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(Revised on June 30, 2015)

Contents

Paper 1

Paper 2

40

Concluding Remarks

80

The year 2013 represented a drastic departure in the history of the JEE.

Till 2012, the selection of the entrants to the IITs was entirely left to the

IITs and for more than half a century they did this through the JEE which

acquired a world wide reputation as one of the most challenging tests for entry

to an engineering programme. Having cleared the JEE was often a passport

for many lucrative positions in all walks of life (many of them having little to

do with engineering). It is no exaggeration to say that the coveted positions

of the IITs was due largely to the JEE system which was renowned not

only for its academic standards, but also its meticulous punctuality and its

unimpeachable integrity.

The picture began to change since 2013. The Ministry of Human Resources decided to have a common examination for not only the IITs, but

all NITs and other engineering colleges who would want to come under its

umbrella. This common test would be conducted by the CBSE. Serious concerns were raised that this would result in a loss of autonomy of the IITs

and eventually of their reputation. Finally a compromise was reached that

the common entrance test conducted by the CBSE would be called the JEE

(Main) and a certain number of top rankers in this examination would have

a chance to appear for another test, to be called JEE (Advanced), which

would be conducted solely by the IITs, exactly as they conducted their JEE

in the past.

So, in effect, the JEE (Advanced) from 2013 took the role of the JEE in the

past except that the candidates appearing for it are selected by a procedure

over which the IITs have no control. So, this arrangement is not quite the

same as the JEE in two tiers which prevailed for a few years. It was hoped

that now that the number of candidates appearing for the JEE (Advanced)

is manageable enough to permit evaluation by humans, the classic practice

of requiring the candidates to give justifications for their answers would be

revived at least from 2014, if not from 2013 (when there might not have

been sufficient time to make the switch-over). But this has not happened

even in 2015 and after a change of regime at the Union Government. The

JEE (Advanced) 2015 is completely multiple choice type and its pattern

differs little from that of the JEE (Advanced) 2014. Questions with single

digit answers come before rather than after those with one or more correct

answers. Also negative marking has been revived for the latter.

Academically (and socially), the JEE (Advanced) has the same status as

the JEE in the past. So, from 2013, the Educative Commentary to JEE

Mathematics Papers is confined only to JEE (Advanced). This year, the

numbering of the questions will be that in Code 8 of the question papers.

As in the past, unless otherwise stated, all the references made are to the

authors book Educative JEE (Mathematics) published by Universities Press,

Hyderabad. The third edition was made available online on the authors blog

http://www.mathjeecommentary.blogspot.in but has now been withdrawn as

the book is easily available in market.

Because of the multiple choice format and many other constraints in papersetting, interesting questions in mathematics are getting rarer. The continuation of these annual commentaries has been possible largely because of

the keen interest shown by the readers. These commentaries are prepared

single-handedly and hence are prone to mistakes of spelling, grammar and occasionally, wrong symbols (but, hopefully, not mistakes of reasoning!). Many

alert readers in the past had pointed out some such errors. They were corrected and the corrected versions were uploaded from time to time. But by

that time the time relevance was reduced.

As an experiment, a draft version of this years commentary on both the

papers was uploaded. Those readers who noticed any errors in it were invited

to send an email to the author at kdjoshi314@gmail.com or send an SMS

to the author at 9819961036. Alternate solutions and any other comments

were also solicited. This really paid off. Several readers, notably Siddhesh

Naik and Deepanshu Rajvanshi pointed out several corrections and suggested

some additions. The elegant solution to Q.55 in Paper 1, where the given

determinant is expanded by writing it as a product of two determinants was

given by Siddhesh Naik. I am grateful to all these and other readers.

2

PAPER 1

Contents

Section - 1 (One Integer Value Correct Type)

13

32

SECTION 1

One Integer Value Correct Type

This section contains eight questions. The answer to each question is a

SINGLE DIGIT INTEGER ranging from 0 to 9, both inclusive.

Marking scheme : +4 If the bubble corresponding to the answer is darkened, 0 In all other cases

Q.41 Let the curve C be the mirror image of the parabola y 2 = 4x with

respect to the line x+y+4 = 0. If A and B are the points of intersection

of C with the line y = 5, then the distance between A and B is

Answer and Comments: 4. Call the line y = 5 as L.

y

L*

C*

A*

x

O

x+y+4=0

B*

C

A

3

Call this parabola C . Then C and C are reflections of each other and

hence are congruent. So C is also a parabola. To find its equation, we

would have to start with a typical point (x0 , y0) on C, find its reflection,

say (x0 , y0), in the line x + y + 4 = 0 and then put y0 2 = 4x0 2 . The

formula for the reflection of a point (x0 , y0) into a line ax + by + c = 0

is possible but rather complicated.

A better approach is to realise that under reflections all distances

are preserved. So, instead of taking the points of intersection of the

parabola C and the line L, we may as well take the points of intersection

of the parabola C and the reflection, say L , of the line L into the line

x + y + 4 = 0. L can be found almost by inspection. The line L cuts

the line x + y + 4 = 0 at the point P = (1, 5) at an angle 45 degrees.

So, L must be the line through P making an angle of 45 degrees with

the line x + y + 4 = 0. Clearly this line is x = 1. This line L happens

to lie along the latus rectum of the parabola C : y 2 = 4x. Hence its

intercept with the parabola has length 4. But even if we miss this,

the points of intersection, say A and B of L with C can be found

by merely solving the two equations y 2 = 4x and x = 1. They are

(1, 2) and so the distance between them is 4. This is also the distance

between A and B since A , B are the reflections of A, B.

This is an excellent problem which tests the ability of realising

how not to do a problem in the most straightforward way and, instead,

look for alternate ways. Once the idea of transforming the problem to

finding the intersections of C and L strikes, the calculation required

is minimal.

Q.42 The minimum number of times a fair coin needs to be tossed, so that

the probability of getting at least two heads is at least 0.96, is

Answer and Comments: 8. Yet another problem where it is much

easier to find the answer by transforming the problem. In the present

case, the transformation is to consider the complementary probability,

say q of the given event. The complementary event here is that at

most one head appears. This falls into two mutually exclusive cases, no

heads and exactly one head. If there are n tosses, then their respective

1

1

1

1

probabilities are ( )n and n ( )n1 = n( )n . Together, q =

2

2

2

2

4

1

(n + 1)( )n . The problem asks for the least integral value of n for

2

which q < 0.04. This reduces to

25(n + 1) < 2n

(1)

As n grows, 2n grows much more rapidly than n + 1. But the least such

n has to be found by trial and error. The L.H.S. is at least 50 and so

n has to be at least 6. By trial we get that 8 is the least integer for

which (1) holds.

This is also a good problem. But the technique of complementary probability is fairly common as compared with using properties of

reflections in the last problem.

Q.43 Let n be the number of ways in which 5 boys and 5 girls can stand in a

queue in such a way that all the girls stand consecutively in the queue.

Let m be the number of ways in which 5 boys and 5 girls can stand in

a queue in such a way that exactly four girls stand consecutively in the

m

queue. Then the value of

is

n

Answer and Comments: 5. To find n think of the 5 girls lined

together as a single object. Then the number of ways to arrange this

object along with the 5 other objects (the boys) is 6!. But the 5 girls

can form a single object in 5! ways. So

n = 6!5!

(1)

row. This object can be formed in 5 4! = 5! ways. Now we have

7 objects, this object with four girls, the remaining girl and the 5

boys. They can be arranged in 7! ways. Hence the total number of

arrangements in which at least 4 girls are together is 7!. But we have

to exclude those in which all 5 girls are together. That is already

counted as n. Moreover, each such arrangement has to be excluded

twice because the excluded girl can be the one at the head or at the

tail. [(G1 G2 G3 G4 G5 ) gets excluded twice, once as G1 (G2 G3 G4 G5 ), and

then again as (G1 G2 G3 G4 )G5 .] Hence

m = 7! 5! 2n = 7! 5! 2 6! 5!

5

(2)

Hence

m

7! 2 6!

=

= 72 =5

n

6!

(3)

problem (by means of forming a single object of the girls) and of complementary counting are present here too. But what makes the problem

tricky is that each unwanted arrangement appears twice. Those who

7! 6!

= 6 as the answer. And since this is also an

miss this will get

6!

integer between 0 and 9, there is no built in alert.

There is an alternate and a less tricky way to calculate m. As

before, form two groups of girls, one with 4 girls and the other with

just one girl. Now arrange the 5 boys on a row in any of the 5! ways.

In each such arrangement, there are six possible places where the two

groups of girls can go, one before the first boy, one after the last boy

and four in between two consecutive boys. We have to insert the two

groups of girls into these 6 places so that they do not go into the same

place. This can be done in 6 5 = 30 ways. This would give

m = 5! 5! 30

(4)

Q.44 If the normals of the parabola y 2 = 4x drawn at the end points of its

latus rectum are tangents to the circle (x 3)2 + (y + 2)2 = r 2 , then

the value of r 2 is

Answer and Comments: 2. Clearly r is the distance of the point

(3, 2) (the centre of the circle) from either of the two normals. The

problem is an unnecessarily clumsy way of asking the distance of this

point from these normals. Actually, either one of the two ends of the

latus rectum will suffice since we are given that both the normals at

both the ends are equidistant from (3, 2). We choose the end A =

dy

(1, 2) for this purpose. The slope of the tangent at this point is

dx

2

dy

2

= = 1 at A. Hence the slope of the

where y = 4x. This gives

dx

y

6

i.e.

x+y3 =0

(1)

|3 2 3|

= 2. Hence r 2 = 2.

1+1

An extremely straightforward problem. A disappointment on the

backdrop of the last three problems. Such problems hardly belong to

an advanced test. Some students may be unnecessarily tempted to use

formulas for equations of normals in terms of their slopes. The direct

approach above is far better.

The distance r of the point (3, 2) is

[x],

0,

x2

x>2

then the value of (4I 1) is

xf (x2 )

dx,

2 + f (x + 1)

x = 2. Since the integrand involves f (x2 ) and f (x + 1), we have to

keep track of when x2 exceeds 2 and also when x + 1 exceeds 2 in the

and the latter at x = 1. Moreover [x] has discontinuities at 0 and 1.

So we have to split

of integration into four subintervals,

the interval

[1, 0], [0, 1], [1, 2] and [ 2, 2] and integrate over each one of them

and add the four integrals. We do so one-by-one. Call the integrand

as g(x) in all cases and denote the four integrals by I1 , I2 , I3 and I4

respectively.

On both the intervals [1, 0] and [0, 1] 0 < x2 < 1 except possibly

at the end-points and so f (x2 ) = [x2 ] = 0. Hence the first two integrals

are 0. In the third and the fourth intervals, x + 1 exceeds 2 except at

the point 1 and hence f (x + 1) = 0. So in both the cases the integrand

x

1

g(x) simplifies to xf (x2 ). On [1, 2] this becomes while on [ 2, 2],

2

2

Z 2

x2 2 1

x

dx = = .

it vanishes. So the whole integral I is merely

2

4 1

4

1

Therefore 4I 1 = 0.

7

the part of the candidate. The integration part itself is trivial.

Q.46 A cylindrical container is to be made from certain solid material with

the following constraints : It has a fixed inner volume of V mm3 , has a 2

mm thick solid wall and is open at the top. The bottom of the container

is a solid circular disc of thickness 2 mm and is of radius equal to the

outer radius of the container. If the volume of the material used to

make the container is minimum when the inner radius of the container

V

is 10 mm, then the value of

is

250

Answer and Comments: 4. A typical problem about the minimisation of a function of one variable. The choice of this variable is ours.

We could take it to be the either the inner or the outer radius of the

base or even its inner or outer height. But since the last part of the

data is in terms of the inner radius of the base, that is a more natural

choice. So let r be the inner radius of the base and h the inner height

of the cylinder. Then its (inner) volume V is

V

= r 2 h

(1)

terms of r. The volume, say W , of the material used is the difference

of the volumes of two coaxial cylinders, the inner one of radius r and

height h and the outer one of radius r + 2 and height h + 2 (and not

h + 4 as the container has no top). Therefore

W = (r + 2)2 (h + 2) r 2 h

= [(r 2 + 4r + 4)(h + 2) r 2 h]

= [h(4r + 4) + 2(r + 2)2 ]

V

= 2 (4r + 4) + 2(r + 2)2

r

(2)

by (1).

As V is a constant, this expresses W as a function of the single variable

dW

must vanish

r. We are given that it is minimum when r = 10. So

dr

at r = 10. By a direct calculation,

dW

4V

8V

= 2 3 + 4(r + 2)

dr

r

r

8

(3)

4V

8V

+

= 48

100 1000

(4)

48V

V

= 48 and gives

= 4.

1000

250

Although this is a problem about minimising a function of one

variable, it is of a different spirit. Here we are not asked to minimise

W . Rather, we are given where it is minimum. That simplifies the

work. If we were asked to find where W is minimum, we would have

to solve a cubic equation and then hunt for the minimum among the

critical points. That makes the problem reasonable. Also the fact that

the answer is an integer between 0 and 9 serves as an alert in case there

are any computational mistakes. This is a good feature of the problem.

(A problem of a similar spirit was asked in Jee 2013 Advanced Paper

1. An open box was to be formed by folding a rectangle after removing

squares of the same size from its four corners and we were given the side

of the square for which this volume is maximum. There is, of course,

nothing wrong in asking problems similar to those in the past years.

But one wishes that the repetition would not have occurred so soon.

Many but not all students must have studied the 2013 papers and those

who did would have an easier time in understanding the problem than

the others. It is all right to repeat an idea that was used last year, or

one that was used a decade ago. But a gap of two years can be unfair.)

which simplifies to

x2 +/6

the region bounded by x = 0, y = 0, y = f (x) and x = a, then f (0) is

Answer and Comments: 3. The problem asks the value of the

function f (x) at x = 0. But this function is not given explicitly. Instead, we are given that f is defined on [0, 1/2] and takes only nonnegative values. We are also told something about the area bounded

by x = 0, y = 0, x = a and y = f (x). This is precisely the area under

the graph of y = f (x). So, the second piece of information means that

Z

a

0

f (x) dx = F (a) + 2

9

(1)

where the R.H.S. needs to be calculated from the first part of the

problem. (We have to assume that (1) holds for all a [0, 1/2]. This

is not clearly stated in the problem. But the problem cannot be solved

without such an assumption.)

Our interest is in f (x). We can get it from (1) by differentiating both

the sides (w.r.t. a) using the second form of the FTC (Fundamental

Theorem of Calculus). We then get,

f (a) = F (a)

(2)

for all a [0, 1/2]. So our problem now reduces to find F (0).

F (x) is defined by an integral in which both the upper and the lower

limits are functions of x. Therefore, by the generalised form of the

second fundamental theorem of calculus,

d 2

d

(x + /6) 2 cos2 x (x)

dx

dx

2 2

2

= 4x cos (x + /6) 2 cos x

(3)

mechanically, let us observe that our interest is only in F (0). So

there is no need to consider those terms in the derivative which are

sure to vanish at 0. The derivative of cos2 x is one such term since it

will involve a factor sin x. As for the derivative of the first term, viz.

4x cos2 (x2 + /6), when we apply the product rule, the factor 4x will

vanish at x = 0. So there is no need to take the derivative of the second

factor cos2 (x2 + /6). The only term in the derivative that remains to

be considered is 4 cos2 (x2

+ /6), evaluated at x = 0. This comes out

3

to be 3 since cos(/6) =

.

2

Problems based on the second form of the FTC are fairly common

in JEE. The present problem is more a test of a candidates ability

to analyse a problem correctly and focus on the essence so as to weed

out unnecessary work. In this sense it is a very good problem. Many

candidates will be tempted to simplify the integrand of F (x) to 1 +

cos 2t. But such a simplification has no role in the solution. In this

respect, the problem is a bit tricky too.

10

Note that the second term (viz. 2) on the R.H.S. of (1) plays no

role in the solution since we are only dealing with the derivative of the

R.H.S. It would thus appear that in the data of the problem, the area of

the region could as well have been given as F (a) + k for any constant

k. But then there would have been an inconsistency. From (3), we

know that F (0) = 2. On the other hand, by putting a = 0 in (1), we

would have gotten F (0) = k. So, the data would be consistent only

if k = 2. A somewhat similar inconsistency in the data had occurred

in a problem in JEE 2011 Paper 1 and was commented upon. (See

Q.22 of that years

Z commentary.) It was given in the statement of the

problem that 6

x Z= 1). The lapse could have been corrected by changing the data to

x

f (t)dt = 3xf (x) x3 5. This would not affect the rest of the

6

1

problem or its solution. Indeed that is probably what led to the lapse.

This year the paper-setters have been careful. The problem could

have been made a little more interesting by giving in the data that the

area of the region was F (a) + k for some constant k and then asking

the candidates to find the value of f (0) + k. That would have forced

the candidates to determine k as 2 from (1) and (3). The answer to

the problem would have been 5 instead of 3, still a single digit number.

Q.48 The number of distinct solutions of the equation

5

cos2 2x + cos4 x + sin4 x + cos6 x + sin6 x = 2

4

in the interval [0, 2] is

Answer and Comments: 8. Superficially, this is a problem of solving

a trigonometric equation. But there is no way to do this unless we first

simplify the expression, say E on the L.H.S. The key idea is to note

that all terms are expressible in terms of sin2 x and cos2 x. So the good

old identity sin2 x + cos2 x = 1 may be useful. Indeed, if we square this,

we get

sin4 x + cos4 x = 1 2 sin2 x cos2 x

11

(1)

sin6 x + cos6 x = 1 3 sin4 x cos2 x 3 sin2 x cos4 x

= 1 3 sin2 x cos2 x(sin2 x + cos2 x)

= 1 3 sin2 x cos2 x

(2)

equation becomes

E=

5

cos2 2x + 2 5 sin2 x cos2 x

4

(3)

cos2 2x 4 sin2 x cos2 x = 0

(4)

cos2 2x sin2 2x = 0

(5)

cos 4x = 0

(6)

or equivalently,

which has eight solutions in the interval [0, 2]. Specifically, the solutions are of the form where 4x = /2 and where 4x = 3/2. But the

problem only asks for the number of solutions.

This is a fairly easy problem once the idea of taking the powers of

the basic identity for sin2 x + cos2 x strikes. See Comment No. 14 of

Chapter 7 for a problem where a similar trick is used.

12

SECTION 2

This section contains TEN questions.

Each question has FOUR options (A), (B), (C) and (D). ONE OR MORE

THAN ONE of these four option (s) is (are) correct.

Marking scheme : +4 If only the bubble(s) corresponding to all the correct

option(s) is (are) darkened, 0 if none of the bubbles is darkened and 2 in

all other cases.

Q.49 Let y(x) be a solution of the differential equation (1 + ex )y + yex = 1.

If y(0) = 2, then which of the following statements is (are) true?

(A) y(4) = 0

(B) y(2) = 0

(C) y(x) has a critical point in the interval (1, 0)

(D) y(x) has no critical point in the interval (1, 0).

Answer and Comments: (A), (C). This is an extremely standard

problem of solving a first order linear differential equation. Normally,

one would begin by recasting the equation in the standard form y +

p(x)x = q(x) and find an integrating factor. In the present problem,

that is hardly necessary. The equation is exact as it stands, because

the L.H.S. is simply the derivative of y(1 + ex ). So, integrating both

the sides, the general solution is

y(1 + ex ) = c + x

(1)

4 = c. So the function y is given by

y(x) =

4+x

1 + ex

(2)

dy

derivative

.

dx

dy

(1 + ex ) (4 + x)ex

=

dx

(ex + 1)2

13

(3)

At a critical point the numerator must vanish. This gives the equation

3ex + xex 1 = 0

(4)

asked is if it has a solution in the interval (1, 0). For this we apply the

Intermediate Value Property. Call the L.H.S. as g(x). It is continuous

2

everywhere. Also g(1) = 1 < 0 since e > 2. On the other hand

e

g(0) = 2 > 0. So, by the Intermediate Value Property, g(x) has at

least one root in (1, 0). Therefore y(x) has at least one critical point

in (1, 0).

The problem is a combination of two unrelated parts. First, solving

an initial value problem and secondly testing if a given function has a

root in an interval. Both are very standard. There is little point in

asking such questions in an advanced test.

Q.50 Consider the family of all circles whose centers lie on the straight line

y = x. If this family of circles is represented by the differential equation

dy

P y +Qy +1 = 0, where P, Q are functions of x, y and y (here y =

dx

2

d

y

and y = 2 ), then which of the following statements is (are) true?

dx

(A) P = y + x

(B) P = y x

2

(C) P + Q = 1 x + y + y + (y ) (D) P Q = x + y y (y )2

Answer and Comments: (B), (C). This problem is about finding

the differential equation of a family of curves. In the present case, a

typical member of the given family is a circle of the form

(x h)2 + (y h)2 = r 2

(1)

family and the differential equation representing it will be of order 2.

To get it we differentiate (1) to get

14

(2)

h=

x + yy

y + 1

(3)

(y + 1)(1 + yy + y 2 ) (x + yy )y

=0

(y + 1)2

(4)

is

(y + 1)(1 + yy + y 2 ) (x + yy )y = 0

(5)

the statement of the problem, the coefficients P and Q are not allowed

to contain y . So we have to recast this equation collecting all the

terms in y together. That gives

(y x)y + (1 + y + y 2 )y + 1 = 0

(6)

the form given in the statement of the problem,

P = yx

and Q = 1 + y + y 2

(7)

(8)

A very mechanical problem. The solution essentially ends at (5).

The remainder is a useless addendum. It might have served some purpose if there were any terms involving y y because then the candidate

would have to think whether to include them as multiples of y or of

y . But since these terms get cancelled, this fine thinking is not tested

anyway.

Q.51 Let g : IR IR be a differentiable

function with g(0) = 0, g (0) = 0

x

g(x), x 6= 0

|x|

and g (1) 6= 0. Let f (x) =

and h(x) = e|x| for all

0,

x=0

15

Then which of the following is (are) true?

(A) f is differentiable at x = 0

(B) h is differentiable at x = 0

(C) f h is differentiable at x = 0 (D) h f is differentiable at x = 0

Answer and Comments: (A), (D). This is a question about the

differentiability of two functions and also of their composites. The

composite of two differentiable functions (when defined) is always differentiable. But sometimes the composite may be differentiable even

when one of the functions fails to be so. An extreme counterexample

is when one of the functions is constant. Then the composite is also

a constant and differentiable regardless of the other function. So such

cases have to be handled carefully.

Let us begin with the differentiability of the function h(x) = e|x|

at x = 0. This is the composite of the absolute value function and the

exponential function. The former is not differentiable at 0. But that is

no reason to hastily declare that the composite e|x| is not differentiable

at 0 because as we just saw, the composite of a differentiable function

with a non-differentiable one can be differentiable sometimes. But we

can put the non-differentiability of |x| to use as follows. We want to

e|x| e0

consider whether lim

exists. For x 6= 0 we rewrite this ratio

x0

x

as

e|x| 1 |x|

e|x| e0

=

x

|x|

x

(1)

derivative of ex at x = 0. But the second factor is 1 for x > 0 and 1

for x < 0. So the product tends to 1 as x 0+ and to 1 as x 0 .

Therefore the product of the two ratios tends to 1 as x 0+ and to 1

as x 0 . Hence h is not differentiable at 0. (There is a slicker way

to see this by observing that |x| = ln(h(x)). The logarithm function is

differentiable everywhere in its domain. So, if h were differentiable at

0, then this composite |x| would be differentiable at 0, a contradiction.)

16

f (x) =

g(x), x > 0

0,

x=0

g(x), x < 0

(2)

g+

(0) which is 0 since it is given that g (0) = 0. Similarly, f (0) will

equal g

(0) which is also 0. Hence f is differentiable at 0. (It would

not be so if g (0) were non-zero.)

We now tackle the differentiability of the composite functions f h

and h f at 0. Note that h(x) is always positive and so by (2)

(f h)(x) = f (h(x)) = g(h(x))

(3)

(h(0)), i.e. g+

(1) = g (1) as g is given to

(f h)+ (0) would equal g+

(1) = g (1).

Since g (1) 6= 0, these two numbers are unequal and so f h is not

differentiable at 0.

Finally, we consider the differentiability of h f at 0. By definition,

(h f )(x) = e|f (x)|

(4)

for all x. Also, by (2), |f (x)| = |g(x)| for all x. Hence e|f (x)| is the

same as e|g(x)| . The exponential function is differentiable everywhere.

So if we can show that |g(x)| is differentiable at 0, then it would follow

that e|g(x)| , and hence h f is differentiable. We are given that g is

differentiable at 0. This does not by itself imply that |g| is differentiable

at x, as one sees from the fact that x is differentiable at 0 but |x| is

not. But now we are also given that g(0) = 0 and g (0) = 0. So, for

x 6= 0 we can write

g(x)

|g(x)| |g(0)|

=

x

x

(5)

with the + sign holding for x > 0 and the minus sign holding for

g(x)

x < 0. But as we are given that

0 as x 0, it follows

x

17

that

g(x)

x

differentiable at 0.

This is a very good conceptual question which demands very fine

thinking and little computations from the candidate. Unfortunately

such questions are unsuitable as multiple choice questions because there

is no way to test if the candidate has understood the subtlety of the

problem, even if he has ticked the correct answer. For example, a

candidate who rules out (B) on the superficial ground that |x| is not

differentiable at 0 cannot be distinguished from one who really proves

it as above.

Negative credit is very unfair for such questions. There is no

reason why a candidate who has done the fine thinking for, say (B),

should be penalised for failure to do a similar thinking for (D). Most

candidates who are adept at the JEE strategy will stay away from such

questions and in the time saved, safely bag 12 points by doing some

absolutely routine questions. So, such questions serve little purpose in

the selection when they are clubbed together with a large number of

mediocre questions.

sin

sin x

and g(x) = sin x for all x IR. Let

6

2

2

(f g)(x) denote f (g(x)) and (g f )(x) denote g(f (x)). Then which

of the following is (are) true?

1 1

1 1

(B) Range of f g is ,

(A) Range of f is

,

2 2

2 2

f (x)

(C) lim

=

(D) There is an x IR such that (g f )(x) = 1

x0 g(x)

6

Answer and Comments: (A), (B), (C). Yet another question involving the composites of two functions. In both the questions the

paper-setters have been careful enough to indicate which of the two

possible interpretations of a composite they have in mind. Although

the interpretation given here is more standard, candidates who follow

the other interpretation should not suffer solely for that reason.

Now, coming to the question itself, all parts except (C) are based on

the images of various intervals under the sine function. and the fact that

18

the range of the composite is the image of the range of the first function.

Since sin x maps IR onto [1, 1], sin x maps it onto [ , ]. The sine

2

2 2

function maps this interval onto [1, 1] and so sin( sin x)) has range

6

2

[ , ]. Under the sine function this interval goes to [1/2, 1/2] which

6 6

proves (A). In (B), we first have to find the range of g which is [ , ].

2 2

Under the sine function, this interval has the same image viz. [1, 1]

as the image of the sine function over the entire IR. So f g has the

same range as f , which is [1/2, 1/2]. In (D), the range of g f is the

image under g of the range of the f which we already know to be the

interval [1/2, 1/2]. So (D) will be true if and only if there is some

2

x [1/2, 1/2] for which g(x) = 1, i.e. sin x = . As the sine function

upon which of the two numbers sin(1/2) and 2/ is bigger. To do this

without calculators, we use the inequality sin x < x for all x > 0. In

particular, sin 1/2 < 1/2. But 1/2 < 2/ since < 4. So (D) is false.

Part (C) is of a totally different spirit than the others. The limit

in question can be calculated rather mechanically using the LHopitals

rule. But a better way out is to put u = g(x) = sin x and v = 6 sin u.

2

Then u 0 and v 0 as x 0. Therefore,

sin( 6 sin u)

f (x)

=

g(x)

u

sin u

sin( 6 sin u)

6

=

sin u

u

6

sin v sin u

=

v

6

u

(1)

.

6

A simple but highly repetitious problem based on the range of the

sine function. Part (D) requires the approximate value of . Part (C)

is a useless addendum.

Both the first and the last factors tend to 1 and so the limit is

Q.53 Let P QR be

a triangle. Let ~a = QR, ~b = RP and ~c = P Q. If

|~a| = 12, |~b| = 4 3 and ~b ~c = 24, then which of the following is (are)

true?

19

|~c|2

|~c|2

|~a| = 12

(B)

+ |~a| = 30

2

2

(A)

computations involving vectors, their lengths and various products.

The lengths of the two sides of the triangle are given. If we were given

the dot product of the vectors representing them, viz. ~a and ~b, we

would also know the angle between them and then we would know the

triangle completely. But we are not given that. We are given |~a|, |~b|,

but not ~a ~b. Instead, we are given ~b ~c. We are also given |~b|. So, if

we could get |~b + ~c|, we would know |~c| and that would also determine

the triangle completely.

The basic idea is that

since the vectors ~a, ~b and ~c

are the sides of the triangle

(directed appropriately), their

vector sum is 0. That is,

~a + ~b + ~c = ~0

(1)

c /2

/2

b

30

30

R

6

6 a O

Therefore

~b + ~c = ~a

(2)

|~b|2 + |~c|2 + 2(~b ~c) = |~a|2 = 144

(3)

This gives

144 48 48 = 48 = 4 3

(4)

and (B).

|~c| =

For the other two statements, we need to know more about the

triangle. Since |~b| = |~c| the triangle P QR is isosceles with P Q = P R.

20

Let the angle between these sides be . Then is the angle between

the vectors ~b and ~c. So,

cos =

~b ~c

1

24

=

=

48

2

|~b| |~c|

(5)

angles are 30 each. Since the angle between the vectors ~b and ~a is

30 the angle between the vectors ~a and ~b is 150 . So we get

3

~a ~b = |~a| |~b| cos 150 = 12 4 3

= 72

(6)

2

which shows that (D) is true. Finlly, for (C), we note that by (1)

~b = ~a ~c. Since ~a ~a = ~0, and the angle between ~a and ~c is also

150, we have

|~a ~b + ~c ~a| =

=

=

=

~a (~a ~c) + ~c ~a

|2~c ~a|

2|~c| |~a| sin 150

4 3 12 = 48 3

(7)

Our solution is purely geometric. The second half of it could have

been shortened a little by resolving the vectors along a suitable pair

of mutually orthogonal unit vectors. We take the midpoint of the side

QR as the origin O and two unit vectors ~i and ~j along OR and OP

1

respectively. Since 6 OP Q = 6 OP R = 60 we have OP = P Q =

2

~

2 3. This allows us to express all the three vectors ~a, b and ~c as linear

combinations of the mutually orthogonal vectors ~i and ~j as

~a = 12~i

~b = 6~i + 2 3 ~j

and ~c = 6~i 2 3 ~j

(8)

(9)

(10)

By a direct calculation,

|~a ~b + ~c ~a| = |~a (~b ~c)|

= |12~i 4 3 ~j| = 48 3

21

(11)

(12)

In fact, the entire solution could have been algebraic. We take ~i to

be a unit vector along ~a and ~j to be a unit vector perpendicular to ~i.

Then, ~a = 12~i and ~b = ~i + ~j where, are some scalars. Then by

(1), ~c = (12 )~i ~j. Since ~b = 4 3, we have

(13)

2 + 2 = (4 3)2 = 48

Similarly, ~b ~c = 24 gives

( + 12) 2 = 24

Solving this system simultaneously, we get

= 6, = 2 3

(14)

(15)

This way we get (9) and (10) more efficiently. Having known all three

vectors in terms of ~i and ~j, all the four statements can be tested one

by one. Thus we see that in the present problem the purely algebraic

solution is fastest. But the gain is not so significant as the problem

itself is simple.

The problem is simple, once the essential idea, viz. Equation (1)

strikes. Unfortunately, there is too much numerical work. Even a single

mistake is costly. So this problem is more a test of speed and numerical

accuracy than reasoning.

Q.54 Let X and Y be two arbitrary, 33, non-zero skew-symmetric matrices

and Z be an arbitrary 3 3 non-zero symmetric matrix. Then which

of the following matrices is (are) skew-symmetric?

(A) Y 3 Z 4 Z 4 Y 3

(C) X 4 Z 3 Z 3 X 4

(B) X 44 + Y 44

(D) X 23 + Y 23

Answer and Comments: (C), (D). Parts (B) and (D) are based

on some simple properties of skew-symmetric and symmetric matrices.

22

matrices is symmetric (respectively, skew-symmetric) and (ii) all powers of symmetric matrices are symmetric (iii) all even powers of skewsymmetric matrices are symmetric while all odd powers of them are

skew-symmetric. With these properties, we have X 44 , Y 44 and hence

their sum symmetric and similarly X 23 + Y 23 skew-symmetric. Note

the analogy of these properties with the properties about the signs of

real numbers. For example, the sum of two positive (negative) real

numbers is also positive (negative). Also the even powers of a negative

number are positive while the odd powers are negative.

However, in (A) and (C), we are dealing with the products of

symmetric and a skew-symmetric matrices. Here the analogy with the

real numbers breaks down. For example, the product of two negative

numbers is positive. But little can be said about the product of two

skew-symmetric matrices unless they commute with each other. So,

(A) and (C) have be handled by directly taking transposes and using

the elementary properties of transposes, viz. the anticommutativity

and the self-reciprocity. (In simpler terms, this means (AB)T = B T AT

and (AT )T = A, for all A, B.)

With these rules in mind, we have

(Y 3 Z 4 Z 4 Y 3 )T = (Y 3 Z 4 )T (Z 4 Y 3 )T

= (Z 4 )T (Y 3 )T (Y 3 )T (Z 4 )T

= Z 4 Y 3 + Y 3 Z 4

(1)

is symmetric. If it were to be skew-symmetric too, it would have to

vanish, which means that Y 3 and Z 4 must commute with each other.

As this is not given, we discard (A).

By an analogous computation,

(X 4 Z 3 Z 3 X 4 )T = (Z 3 )T (X 4 )T (X 4 )T (Z 3 )T = Z 3 X 4 X 4 Z 3

since Z 3 and X 4 are both symmetric. This shows that X 4 Z 3 Z 3 X 4

is skew-symmetric.

Because of their limited scope, questions about matrices tend to be

repetitious. The present one is a little unusual. The calculations are

23

possibility that the matrix may be skew-symmetric as well, as pointed

out above. It is given in the question that the matrices Y, Z are nonzero. But that does not preclude the possibility that Z 4 commutes

with Y 3 . For example, Z could be the identity matrix. It would have

been better if the question had asked to identify those matrices that

are necessarily skew-symmetric.

Q.55 Which of the following values of satisfy the equation

(1 + )2 (1 + 2)2 (1 + 3)2

(2 + )2 (2 + 2)2 (2 + 3)2

(3 + )2 (3 + 2)2 (3 + 3)2

(A) 4 (B) 9

(C) 9

= 648 ?

(D) 4

is a polynomial in . It would be horrendous to compute D by direct

expansion. But if we subtract R2 from R3 and then R1 from R2 we get

1 + 2 + 2 1 + 4 + 42 1 + 6 + 92

3 + 2

3 + 4

3 + 6

D=

5 + 2

5 + 4

5 + 6

(1)

Next, we subtract the middle row from the other two to get

2 2 42 2 92 2

D = 3 + 2 3 + 4 3 + 6

2

2

2

(2)

2

42

92

D = 3 + 2 3 + 4 3 + 6

2

2

2

(3)

To simplify D further, we subtract the first column from the other two

to get

2

32 82

D = 3 + 2 2 4

2

0

0

24

(4)

D = 2(123 163) = 83

(5)

roots are 0, 9 and 9.

evaluation of determinants by various manipulations. But they used to

be full length questions, allowing five to ten minutes. In the present setup, there are hardly 3 minutes. A saving feature is that the candidates

do not have to explain their work. In the present problem, it is unlikely

that a candidate would get the correct answer by mere substitution. So

this is a good, classic type problem on determinants.

There is an alternate, albeit a trickier solution to the problem.

We first expand all the squares and then write the determinant as the

product of two determinants, viz.

D =

1 + 2 + 2 1 + 4 + 42 1 + 6 + 92

4 + 4 + 2 4 + 8 + 42 4 + 12 + 92

9 + 6 + 2 9 + 12 + 42 9 + 18 + 92

1

1

1 1 1 1

4 2 1 2 4 6

9 3 1 2 42 92

3

2

1 1 1 1 1 1

4 2 1 1 2 3

9 3 1 1 4 9

(6)

last column from the others for the first and subtracting the first row

from the remaining ones for the second. Their values are 2 and 2.

(It is not an accident that their values are the negatives of each other,

because if we interchange the first and the last column of the first determinant and then take its transpose, we get the second determinant.)

So, we finally get D = 83 which is the same as (5). The rest of the

solution is the same.

More generally one can consider a determinant D of the form

D = D(x, y, z, a, b, c) =

25

(b + x)2 (b + y)2 (b + z)2

(c + x)2 (c + y)2 (c + z)2

(7)

where x, y, z, a, b, c are variables. When fully expanded this is a homogeneous polynomial in these six variable in which each term has total

degree 6. By expanding the squares and taking steps similar to those

above, one can show that

D(x, y, z, a, b, c) =

a2 a 1 1 1 1

b2 b 1 x y z

c2 c 1 x2 y 2 z 2

(8)

(3.26)) and hence has value (y x)(z y)(z x). If we interchange

the first and the last columns of the first determinant and take its

transpose, that is also a Vandermonde determinant. As a result, we

get

D = 2(a b)(b c)(c a)(x y)(y z)(z x)

(9)

There is also an easier way to see this if we observe from (7) that D

vanishes if, say a = b. Hence (a b) is a factor of D. And so are

(b c), (c a), (x y), (y z) and (z x). So the product of these

six factors also divides D(x, y, z, a, b, c). But this product is already a

polynomial of total degree 6 in x, y, z, z, b, c. Hence we must have

D(x, y, z, a, b, c) = k(a b)(b c)(c a)(x y)(y z)(z x) (10)

for some constant k. The value of k can be determined by giving some

special, simple values to a, b, c, x, y, z, e.g. a = x = 0, b = y = 1 and

c = z = 1. (A similar technique is also possible for evaluating the

Vandermonde determinant.)

There is a certain formal resemblance between the determinant in

(7) and the determinant

D(P, Q, R, A, B, C) = cos(B P ) cos(B Q) cos(B R)

cos(C P ) cos(C Q) cos(C R)

(11)

show that this determinant vanishes for all A, B, C, P, Q, R. Again, a

direct expansion is simply ruled out. But a solution is possible using

26

elementary row operations and trigonometric identities. The best solution, however, is to expand all the entries of the determinant and then

show that it equals the product

cos A sin A 0

cos B sin B 0

cos C sin C 0

sin P sin Q sin R

0

0

0

(12)

plane, different from P1 and P2 which passes through the intersection

of P1 and P2 . If the distance of the point (0, 1, 0) from P3 is 1 and the

distance of a point (, , ) from P3 is 2, then which of the following

relations is (are) true ?

(A) 2 + + 2 + 2 = 0 (B) 2 + 2 + 4 = 0

(C) 2 + 2 10 = 0 (D) 2 + 2 8 = 0

Answer and Comments: (B), (D). This is a problem about a

parametrised family of planes. Let L be the line of intersection of

the planes P1 and P2 . Write the equations of P1 and P2 in the form

E1 = 0 and E2 = 0 where E1 , E2 are linear expressions in x, y, z.

Then the equation of every plane passing through L is of the form

E1 + E2 = 0 for some values of the parameters and . We can

dispense with one of the parameters, say and consider an equation of

the form E1 + E2 = 0. This will represent all possible planes through

L for various values of , except the plane P1 . (For the plane P1 ,

we need = 1 and = 0.) Similarly, the equation E1 + E2 = 0

will represent all planes through L except P2 . In the present case we

are given that P3 is different from both P1 and P2 . So we are free to

take either approach. (See Comment No. 13 of Chapter 9) for more

examples of this technique.)

We take E1 as y, E2 as x + z 1. Then the equation of P3 is of the

form

x + z 1 + y = 0

(1)

for some value of . To determine it, we use the condition that the

distance of the point (0, 1, 0) is 1. This gives

0 + + 0 1

1 + 2 + 1

27

=1

(2)

which simlifies to

( 1)2 = 2 + 2

(3)

1

and determines as . So the equation of the plane P3 is x + z

2

1

1 y = 0 or, equivalently,

2

2x y + 2z 2 = 0

(4)

1. This means

This means

2 + + 2 2

4+1+4

=2

(5)

2 + 2 2 = 6

(6)

A routine problem once the idea of a parametrised family strikes.

A discerning student will observe that normally, (2) would reduce to a

quadratic in , which is consistent with the fact in the family of planes

containing the line L, there are two planes whose distance from the

point (0, 1, 0) is 1. In the present case, one of these two planes is P1

itself and is discarded by the data. Since we took the equation of P3

in the form E1 + E2 = 0, which represents all planes containing L

except P1 , (2) degenerated into a linear equation in . Had we taken

the equation of P3 in the form

y + (x + z 1) = 0

(7)

1

1 + 22

=1

(8)

its roots. The root = 0 gives the plane P1 and has to be discarded.

28

The second root, viz. 2 will give the equation of the plane P3 as

2x y + 2z 2 = 0 which is the same as (4). So the answer does

not change if we take the other parametrisation of the family, but the

work involved does slightly. So, we were rather lucky to start with the

equation E1 +E2 = 0 rather than with E1 +E2 = 0. A sharp student

will, however, not leave this choice to luck. He will observe that P1 is

already at a distance 1 from (0, 1, 0) and since it is to be excluded, it

is safe (and numerically easier) to take the equation of P3 as (1) rather

than (7). But in the present problem the advantage gained is minor.

Most candidates would anyway prefer to start with (1) rather than (7)

because it is simpler. If the advantage were substantial and the easier

option not so tempting, then this would have been a good problem

which rewards the sharp candidates.

Q.57 In IR3 , let L be a straight line passing through the origin. Suppose

that all the points on L are at a constant distance from the two planes

P1 : x + 2y z + 1 = 0 and P2 : 2x y + z 1 = 0. Let M be the

locus of the feet of the perpendicular drawn from the points on L to

the plane P1 . Which of the following points lie(s) on M ?

(A) (0, 5/6, 2/3) (B) (1/6, 1/3, 1/6)

(C) (5/6, 0, 1/6)

(D) (1/3, 0, 2/3)

Answer and Comments: (A), (B). The first part of the data simply

means that the line L is parallel to both P1 and P2 and hence to the

line of their intersection. From the two equations of the planes, viz.

x + 2y z = 1

and 2x y + z = 1

(1)

(2)

1, 3, 5. (This is a standard result. For those who dont know it,

these are the components of the vector u v where u = i + 2j k and

v = 2i j + k are normals to the planes P1 , P2 respectively.) Since L

passes through the origin, its parametric equations are

x = t, y = 3t, z = 5t

where t is a parameter.

29

(3)

the points on L to the plane P1 will form a line M parallel to L. To

determine M, it suffices to know any one point P0 = (x0 , y0 , z0 ) on it.

We take it to be the foot of the perpendicular from the point (0, 0, 0)

(which is given to lie on L) to P1 . To determine this point we need

three equations in x0 , y0 , z0 . one of them comes of course, from the

equation of the plane P1 . That is

x0 + 2y0 z0 = 1

(4)

The other two equations come from the fact that the vector OP0 =

x0 i + y0 j + z0 k is perpendicular to P1 and hence parallel to the normal

vector u = i + 2j k. This gives

x0 = r, y0 = 2r, z0 = r

(5)

for some real number r. Substituting this into (4), we get r = 1/6.

Hence P0 = (1/6, 1/3, 1/6). Therefore the locus M is the line

x + 1/6

y + 1/3

z 1/6

=

=

=

1

3

5

(6)

M. For the point (0, 5/6, 2/3) to lie on M, we must have = 1/6

from x + 1/6 = . This value also satisfies the other two equations

in (6). But for the other two given points, we have y = 0, whence

= 1/9. But that would make x = 1/9 1/6 = 5/18. So the

points in (C) and (D) do not lie on M.

A fairly simple problem once the idea strikes that L is parallel to

the line of intersection of the two given planes. In the conventional

examination the solution would end with (6), i.e. finding the equation

of the locus M. Asking which of the given four points satisfy (6) is

sheer arithmetic and prone to numerical errors. Also it is anybodys

guess what is the reason for making the candidates do this work four

times. The only answer is that the constraints on the paper-setters

stipulate that every MCQ must have four choices. Thank God it was

not 10.

Q.58 Let P and Q be distinct points on the parabola y 2 = 2x such that

a circle with P Q as a diameter passes through the vertex O of the

30

parabola. IfP lies in the first quadrant and the area of the triangle

OP Q is 3 2, then which of the following are the coordinates of P ?

Answer and Comments: (A), (D). There is a minor anomaly in the

wording of the problem. Given any two distinct points P and Q, there

is only one circle having P Q as a diameter. So instead of saying a

circle with P Q as a diameter, the wording should have been the circle

with P Q as a diameter. This might be just a lapse on the part of the

paper-setters. But it might confuse a discerning student, maybe only

for a few seconds. But in a severely competitive test even a few seconds

are precious.

Now, coming to the problem itself, take the points P and Q in the

parametric form as

t21

P = ( , t1 )

(1)

2

t2

(2)

and Q = ( 2 , t2 )

2

2

2

and

respectively. As O lies

Then the slopes of OP and OQ are

t1

t2

on the circle with P Q as a diameter, we have OP OQ which means

4

= 1 and hence

t1 t2

t1 t2 = 4

(3)

OP Q is right angled at O, we get

1

OP OQ

3 2 =

2

|t2 | q 2

1 |t1 | q 2

t1 + 4

t2 + 4

=

2

2

2

q

1

=

(t21 + 4)(t22 + 4)

(4)

2

using (3). Squaring both the sides

(t21 + 4)(t22 + 4) = 72

31

(5)

t21 + t22 = 10

Hence (t1 t2 )2 = 10 2t1 t2 = 18. This gives

t1 t2 = 3 2

(6)

(7)

and hence t2 negative by (3) again. So the positive sign must hold.

Now that we know both t1 t2 and t1 t2 , we can find t1 by solving the

equation

t1 +

4

=3 2

t1

(8)

t21 3 2t1 + 4 = 0

(9)

3 2 2

i.e. 2 2 and 2. The corresponding point

The roots are

2

t21

P = ( , t1 ) then is (4, 2 2) and (1, 2).

2

A routine problem, once the idea of taking the points in their

parametric forms strikes.

SECTION 3

This section contains TWO questions. Each question contains two columns

Column I and Column II. Column I has four entries (A), (B), (C) and

(D). Column II has five entries (P), (Q), (R), (S) and (T). Match the entries in Column I with the entries in Column II. One or more entries in

Column I may match with one or more entries in Column II.

Marking scheme: For each entry in Column I, 2 points if fully correct, 0

points if not attempted and 1 points in all other cases.

32

Q.59

Column I

(A) In IR2 , if the magnitude ofthe projection

vector

of the vector

i + j on 3 i + j is 3 and if

= 2 + 3, then possible value(s) of || is

(are)

(B) Let a and b be(real numbers such that the func3ax2 2, x < 1

tion f (x) =

is differenbx + a2 ,

x1

tiable for all x IR, then the possible value(s)

of a are

Column II

(P) 1

(Q) 2

(R) 3

(S) 4

(T) 5

(3 3 + 2 2 )4n+3 + (2 + 3 3 2 )4n+3 + (3 +

2 + 3 2)4n+3 = 0, the possible value(s) of n is

(are)

(D) Let the harmonic mean of two positive real numbers a and b be 4. If q is a positive real number

such that a, 5, q, b is an arithmetic progression,

then the value(s) of |q a| is (are)

Answers and Comments: (A;P,Q), (B;P,Q), (C; P,Q,S,T), (D; Q,T).

The four parts in Column I are quite independent of each other. But

so is the marking for them. This is a merciful departure from the past

where credit would be given only if all parts of Column I were answered

correctly. But then one wonders what was the point of grouping these

questions in such an elaborate manner especially when all of them have

numerical answers. Why not ask four separate questions, each having

five possible answers? This is again a silly constraint on the papersetters regarding the format of the question paper.

Now, coming to the questions in Column

I, in (A), a unit vector in

3

1

i + j. Hence the projection

the direction of the vector 3 i + j is

2

2

3

1

+ |. Equating this with 3 and

of i + j has magnitude |

2

2

squaring gives

( 3 + )2 = 12

(1)

33

2

We are further given that = . Substituting this into (1) we get

3

(4 2)2 = 36

(2)

= 1. Hence || = 1 or 2.

we need two equations in them. One is provided by the fact that

every differentiable function is continuous. Continuity of f (x) at x = 1

implies

3a 2 = b + a2

(3)

As for differentiability of f (x), the right handed and the left handed

derivatives at x = 1 are 6a and b respectively. So we get

b = 6a

(4)

a2 3a + 2 = 0

(5)

In (C), we are dealing with a sum of the (4n + 3)-th powers of three

quadratic expressions in , specifically,

E1 = 3 3 + 2 2

E2 = 2 + 3 3 2

and E3 = 3 + 2 + 3 2

(6)

(7)

(8)

polynomials. This suggests that the polynomials must be related to

each other in some simple way. (If no such relationship exists, then the

problem would be extremely hard.) Once this idea strikes, finding the

actual relationship is easy. Since 3 = 1, 4 = , 5 = 2 we see that

E1 = E3 and E2 = 2E3

34

(9)

( 4n+3 + ( 2 )4n+3 + 1)E34n+3 = 0

(10)

factorising 3 1 and noting that 6= 1.) So E3 = 3 + 2 3 3 =

6 which is non-zero. Hence E34n+3 is also non-zero. So, only the

first factor of the L.H.S. of (10) vanishes. Further, since 3 = 1, 3n

and ( 2 )3n also equal 1 each. So the equation simplifies to

n + 2n + 1 = 0

(11)

equation becomes 1 + 1 + 1 = 0 which is impossible. However, for other

values of n, it reduces to + 2 + 1 = 0 (for n = 1, 4, 7, . . .) or to

2 + + 1 = 0 (for n = 2, 5, 8, . . .) both of which are true. Hence n

must not be divisible by 3. In Column II, the possible values of n are

those other than 3.

Finally, in (D), the first part gives the equation

ab = 2(a + b)

(12)

We are also given that a, 5, q, b are in an A.P. This gives two more

equations, viz.

q + a = 10

and b + 5 = 2q

(13)

(14)

Eliminating q,

b = 2(10 a) 5 = 15 2a

(15)

a(15 2a) = 2a + 30 4a

(16)

2a2 17a + 30 = 0

35

(17)

17 7

, i.e. 6 and 5/2. Corresponding values of q are 4

4

and 15/2. So |q a| = 2 in the first case and |q a| = 5 in the second.

(C) is the only interesting problem in the bunch. The others reduce

to solving systems of two equations in two unknowns. There is nothing

very great either in formulating these equations or in solving them.

These parts hardly belong to an advanced test.

Q.60

Column I

(A) In a triangle XY Z, let a, b and c be the

lengths of the sides opposite to the angles X, Y

and Z respectively. If 2(a2 b2 ) = c2 and

sin(X Y )

=

, then possible values of n for

sin Z

which cos(n) = 0 is (are)

(B) In a triangle XY Z, let a, b and c be the

lengths of the sides opposite to the angles X, Y

and Z, respectively. If 1 + cos 2X cos 2Y =

a

2 sin X sin Y , then possible value(s) of is (are)

b

the position vectors of X, Y and Z with respect

to the origin O, respectively. If the distance of Z

from the bisector of the acute angle of OX with

3

OY is , then possible value(s) of || is (are)

2

Column II

(P) 1

(Q) 2

(R) 3

(S) 5

(T) 6

bounded by x = 0, x = 2, y 2 = 4x and y =

|x 1| + |x 2| + x, where {0, 1}. Then

8

the value(s) of F () +

2 when = 0 and

3

= 1 is (are)

Answers and Comments: (A; P,R,S), (B; P), (C; P,Q), (D; S,T).

In this question, items (A) and (B) have a common setting. Since

a, b, c are proportional to the sines of their opposite angles, the equation

36

2(sin2 X sin2 Y ) = sin2 Z

(1)

2 sin(X Y ) sin(X + Y ) = sin2 Z

(2)

sin(X Y )

1

=

sin Z

2

(3)

1

get = and the equation to be solved reduces to

2

cos(n/2) = 0

(4)

Column II are 1, 3 and 5.

In (B), the condition given is in terms of the angles and we have

a

sin X

which equals

. Let us first recast the given

to find the ratio

b

sin Y

condition in terms of sin X and sin Y .

1 + cos 2X 2 cos 2Y

= (2 2 sin2 X) 2 + 4 sin2 Y

= 4 sin2 Y 2 sin2 X

(5)

2 sin2 Y sin2 X = sin X sin Y

(6)

sin X

as (which is the ratio we are

sin Y

interested in), we get a quadratic in , viz.

Dividing by sin2 Y and calling

2 + 2 = 0

(7)

whose roots are = 1 and 2. As , being the ratio of two sides cannot

be negative, we get = 1.

37

In (C), the vectors are involved only superficially. The first part

ofthe datais just another way of saying that the points X, Y, Z are

( 3, 1), (1, 3) and (, 1 ) respectively. Even without drawing a

diagram it is obvious that the points X and Y are symmetrically located

w.r.t. the line y = x. Also OX and OY are inclined at angles 30 and

60 respectively. So, the line y = x is the acute angle bisector 6 XOY .

3

We are given that the distance of Z from this line y x = 0 is .

2

This implies

3

|(1 ) |

=

1+1

2

(8)

or 1.

Finally, (D) consists of two separate problems, one for = 0 and

the other for = 1. The first one is easier because in this case y = 3

and so, F (0) is the area bounded by x = 0, x = 2, y 2 = 4x and y = 3.

It is the shaded area OABCO in the figure below.

y

C

(0, 3)

y=3

A

x

O

(2, 0)

By a direct calculation,

F (0) =

8 2

4 3/2 2

3 2 x dx = 3x x = 6

0

3

3

(9)

8 2

= 6 which tallies with (T) in Column II.

So, F (0) +

3

On the other hand, when = 1,

y = |x 1| + |x 2| + x

38

(10)

y = |x 1| + 2

(11)

one from C = (0, 3) to A = (1, 2) and the other from A = (1, 2) to

B = (2, 3). So this time F (1) is the sum of the two shaded areas

OACO and ADBA shown in the figure below.

y

C

(0, 3)

y=3

D

x

O

(2, 0)

F (1) =

Z 2

3 x 2 x dx +

x + 1 2 x dx

1

(12)

8 2

.

By a routine integration which we skip, F (1) comes out to be 5

3

8 2

= 5.

So F (1) +

3

Parts (A) to (C) are straightforward, but too elementary to be

asked in an advanced test. The wording of (D) is clumsy and many

candidates might not understand the problem and might have skipped

it. They are the clever ones, because those who do struggle successfully to realise that the problem involves the calculation of two (in

reality three) unrelated areas will pay a heavy price in terms of time

and the strong possibility of some numerical error. If the idea was

merely to give a problem about identifying and finding the area of a

plane region, the first half of the problem (where = 0) would have

39

served the purpose quite well. Adding one more part which hardly

tests anything new (except the ability to draw the graph of the function y = |x 1| + |x 2| + x) is nothing short of torture. It is such

sadistic problems which make success at the JEE a matter of adopting a clever strategy, whose prime rule is to simply stay away from a

problem which is clumsily worded and utilise the time saved on routine

problems requiring mediocre intelligence.

PAPER 2

Contents

Section - 1 (One Integer Value Correct Type)

40

54

71

SECTION 1

One Integer Value Correct Type

This section contains eight questions. The answer to each question is a

SINGLE DIGIT INTEGER ranging from 0 to 9, both inclusive.

Marking scheme : +4 If the bubble corresponding to the answer is darkened, 0 In all other cases

Q.41 Suppose that all the terms of an arithmetic progression (A.P.) are natural numbers. If the ratio of the sum of the first seven terms to the

sum of the first eleven terms is 6 : 11 and the seventh term lies between

130 and 140, then the common difference of this A.P. is

Answer and Comments:

viz. its initial term, say a,

by d. To determine them,

problem gives only one. In

and its common difference, often denoted

we need two equations in a and d. The

such cases some additional restrictions on

40

the variables such as that they are integers or that they lie in some

specific intervals is needed to determine the unknowns. The present

problem is of this type. (See Exercise (4.24) for an example where

there are fewer equations than unknowns and still a unique solution.)

With the notations just introduced, the sum of the first n terms

(often denoted by Sn ) is

Sn = na + d(1 + 2 + . . . + (n 1)) = na +

We are given that

S7

S11

6

.

11

n(n 1)

d

2

(1)

7a + 21d

6

=

11a + 55d

11

(2)

a = 9d

(3)

the second condition. The seventh term is a + 6d. So we are given

130 < a + 6d < 140

(4)

(5)

140 is 135. Setting 135 = 15d gives d = 9.

A simple, but thought provoking problem. The computations

needed are minimal and can be done quickly. Such questions are ideal

for multiple choice tests because the answer is not likely to come by

guessing.

Q.42 The coefficient of x9 in the expansion of (1 + x)(1 + x2 )(1 + x3 ) . . . (1 +

x100 ) is

Answer and Comments: 8. The problem is superficially algebraic.

But in reality it is combinatorial. When the product, say P is fully

41

terms of the form u1 u2 . . . u100 where each ui has two possibilities , either

1 or xi , for i = 1, 2, . . . , 100. We can combine these two possibilities

by saying that ui = xni where the exponent ni is either 0 or i. So,

u1 u2 . . . u100 will equal un1 +n2 +...+n100 . Clearly only the positive terms

in the exponent matter. So each occurrence of x9 will result from a

sequence of distinct positive integers in ascending order adding to 9.

Once this is understood, the problem reduces to finding the number of ways to express 9 as a sum of distinct integers (arranged in an

ascending order), ranging from 1 to 100. As 9 is a small number, these

possibilities can be counted by classifying according to the number of

terms in the summation. There is only one way to express 9 as a sum

of a single integer, viz. 9 = 9. To find the number of ways to write 9

as a sum of two distinct integers in ascending order, we go systematically as 1 + 8, 2 + 7, 3 + 6 and 4 + 5. If there are three summands, say

9 = a + b + c, with 0 < a < b < c, note that a + b is at least 3. So

c can be only 6, 5 or 4 since with c = 3, a < b < c implies a + b can

be at most 3. For c = 4, we have a + b = 5. But since a < b < c this

can happen only a = 2, b = 3. For c = 5, a + b will equal 4 only for

a = 1, b = 3 and finally for c = 6, the only possibility is a = 1, b = 2.

There is no need to go further because the sum of any four distinct

positive integers will be at least 10. So 9 can be expressed as a sum of

distinct positive integers in an ascending order in 1 + 4 + 3 = 8 ways.

The problem could have been formulated as a combinatorial

problem. Suppose there are 100 boxes, numbered 1 to 100. Then the

problem asks for the number of ways to put nine identical balls into

these boxes, so that each box is either empty or contains as many balls

as its number.

In the present case, we have translated the algebraic problem

into a combinatorial one. It could have as well been done the other

way. The work done is essentially the same with either approach. But

there are situations where converting a combinatorial problem into an

algebraic one pays off. Consider, for example, the problem of counting

selections with repetitions. So, let an,k be the number of ways to choose

k objects from n types of objects with repetitions allowed freely. This

is equivalent to placing k identical balls into n distinct boxes, there

42

k+n1

by thinking of

is a tricky way to show that an,k equals

k

each such placement as an arrangement of k balls and n 1 inter-box

partitions. (See Comment No. 5 of Chapter 1.)

The algebraic version of this problem can be constructed as follows.

Each of the n boxes can hold at least 0 and at most k balls. So an,k is

the coefficient of xk in the expansion of (1 + x + x2 + . . . + xk )(1 + x +

. . . + xk ) . . . (1 + x + . . . + xk ), there being n factors in all. Equivalently,

an,k is the coefficient of xk in (1 + x + x2 + . . . + xk )n . This conversion

does not help much by itself. But, instead of taking the polynomial

1 + x + x2 + . . . + xk , let us take the entire infinite series 1 + x + x2 +

. . . + xk + xk+1 + xk+2 + . . .. This may appear useless because the extra

terms we are adding cannot contribute to any selection. (There is no

way to write k as a sum of non-negative integers if one or more terms

is greater than k.) But now the advantage is that the infinite series

1

1 + x+ x2 + . . .+ xk + . . . can be identified as the geometric series

.

1x

Doing

thisfor each of the factors, we see that an,k is the coefficient of xk

n

1

. This too, is not of much use by itself. But if rewrite this

in

1x

expression as (1x)n and expand it using the binomial theorem where

the exponent is the negative integer n, we see that the coefficient of

(n)(n 1) . . . (n k + 1)

= (1)k

xk is (1)k n

which comes

k

k!

n(n + 1) . . . (n + k 1)

which is nothing but the binomial

out to be

k!

!

n+k1

.

coefficient

k

What makes this solution possible is the rich machinery of algebra,

including power series. The combinatorial solution was elementary but

rather tricky. The situation is analogous to the relationship between

pure geometry and coordinate geometry. Pure geometry solutions are

elegant but sometimes tricky. When coordinates are introduced the

powerful machinery of algebra makes the problem more amenable, if

somewhat dull. (See, for example, the second proof of the concurrency

of the three altitudes of a triangle, given in Comment No. 3 of Chapter

8.)

43

recasting of a combinatorial problem in terms of a suitable polynomial

is sometimes useful when the answer is less demanding than finding the

coefficient of a specific power of x. Consider, for example, the problem

of placing n identical balls into, say, r distinct boxes with the restriction

that each box can contain at most k balls. If we denote this number by

an , then as we just saw, an is the coefficient of xn in the polynomial,

say f (x) = (1 + x + x2 + . . . + xk )r and there is no easy of finding

it except for some select values of n. For example, we can say that

ark = 1 and an = 0 for n > rk because f (x) is a polynomial of degree

rk and leading coefficient 1. But this is something obvious by common

sense anyway. Similarly, we can tell a0 = 1 and a1 = r equally easily

with or without the help of the polynomial f (x).

But, suppose that our problem is not to calculate an for a particular

value of n, but to find the entire sum, say S1 = a0 +a1 +a2 +. . .+ark . We

could have written S1 ostensibly as an infinite sum S1 =

an because

n=0

all except finitely many terms of this series vanish. Since f (x) = a0 +

a1 x + a2 x2 + . . . ark xrk , we have S1 = f (1). From the factorisation of

f (x) as (1 + x + x2 + . . . + xk )r we immediately get S1 = (k + 1)r .

Combinatorially, S1 is the number of all possible ways of putting any

numbers of identical balls into r distinct boxes so that each box contains

at most k balls. Here, too, a direct combinatorial argument is easy

because, for each of the r boxes, there are k + 1 possibilities depending

upon how many balls go into it. So again, this example does not quite

bring home the power of algebraic codification.

To do so, consider S2 = a0 + a2 + +a4 + . . . + a2m + . . ., i.e. the sum

of the coefficients of all even degree terms in f (x). Combinatorially, S2

is the number of all possible ways to put any numbers of balls into the

boxes as before with the additional restriction that the total number

of balls is even. (It is not required that only even numbers of balls

go into the individual boxes. The restriction of evenness is only on

the total number of balls.) This time, the combinatorial count is not

as immediate as for S1 . But the algebraic one is simple. If we add

f (1) and f (1), the even powers add up while the odd ones cancel

f (1) + f (1)

out. In other words, S2 = a0 + a2 + a4 + . . . =

. This is

2

44

f (1) = (1 1 + 1 . . . + (1)k+1 )r which equals 0 or 1 according as

(k + 1)r

(k + 1)r + 1

k is odd or even. Hence S2 =

or

depending upon

2

2

whether k is odd or even.

The essence of the calculations of S1 and S2 was, respectively, that

for every non-negative integer n, 1n = 1 for all n while 1n +(1)n equals

2 or 0 depending upon whether n is a multiple of 2 or not. Note that

1 and 1 are the square roots of 1. In the solution to Part (C) of Q.59

of Paper 1, we proved that 1n + n + ( 2 )n equals 3 or 0 depending

upon whether n is a multiple of 3 or not. So, for any polynomial

f (x) = a0 +a1 x+a2 x2 +. . ., the sum, say S3 = a0 +a3 +a6 +. . .+a3m +. . .

f (1) + f () + f ( 2 )

would equal

.

3

The picture is now quite clear. Even if we may not be able to identify

the coefficients an individually for all n, we can calculate the sums of the

coefficients whose suffixes are in an A.P. by adding the values of f (x) at

the d complex d-th roots of unity where d is the common difference of

this progression. Sometimes such sums have some appeal and provide

an unexpected solution to a problem where a direct combinatorial count

may be laborious. Consider, for example, the problem of finding the

number of 6-digit numbers whose digits add to a number of the form

4p + 1. This is equivalent to counting the number of ways to put

identical balls into 6 boxes the first of which can hold 1 to 9 balls and

the remaining 0 to 9 balls each, so that the total number of balls is of

the form 4p + 1. Then the number we want is the sum a1 + a5 + a9 + . . .

for the polynomial f (x) = (x + x2 + . . . + x9 )(1 + x + x2 + . . . + x9 )5 .

Equivalently, this is the sum b0 + b4 + b8 + . . . for the polynomial g(x) =

b0 +b1 x+b2 x2 +. . . = (1+x+x2 +. . .+x8 )(1+x+x2 +. . .+x9 )5 . By our

g(1) + g(i) + g(1) + g(i)

. By direct substitutions

work, the count is

4

of the powers of these complex numbers into the factors of g(x), we have

g(1) = 900000, g(1) = 0, g(i) = (1+i)5 = 44i and g(i) = 4+4i.

9000000 4 4

So the desired number is

= 249998.

4

x2 y 2

Q.43 Suppose that the foci of the ellipse

+

= 1 are (f1 , 0) and (f2 , 0)

9

5

where f1 > 0 and f2 < 0. Let P1 and P2 be two parabolas with a com45

mon vertex at (0, 0) and with foci at (f1 , 0) and (2f2 , 0), respectively.

Let T1 be a tangent to P1 which passes through (2f2 , 0) and T2 be a

tangent to P2 which passes through (f1 , 0). If m1 is!the slope of T1 and

1

m2 is the slope of T2 , then the value of

+ m22 is

m21

Answer and Comments: 4. The problem is about the tangents to

the two parabolas. The role of the ellipse is only to specify

s the foci of

2

5

the two parabolas. The eccentricity e of the given ellipse is 1 =

9

3

and so the foci are at (3e, 0) i.e. f1 = 2 and f2 = 2.

Since both the parabolas P1 and P2 have their vertices at (0, 0), and

their foci are at (2, 0) and (4, 0) respectively, their equations are

y 2 = 8x

y 2 = 16x

(1)

(2)

parabola y 2 = 8x. So it intersects the parabola in two coincident

points. Therefore, the quadratic m1 (x + 4)2 8x = 0 has discriminant

0. On simplification, this gives (8m21 8)2 = 64m41 and hence

m21 =

1

2

(3)

(Incidentally, this shows that there are two lines through (4, 0) that

touch the parabola P1 . The problem does not specify which of the two

is to be chosen. But the answer does not depend on the choice since it

involves only m21 .)

By an entirely analogous computation, the slope m2 of T2 satisfies

m22 = 2

(4)

1

+ m22 equals 4.

m21

Those who know the equation of a tangent to a parabola in terms of

its slope can shorten the work by observing that a tangent to y 2 = 8x

having slope m1 has its equation of the form

y = m1 x +

46

2

m1

(5)

and get (3) from the fact that this tangent passes through (4, 0). Similarly an alternate derivation of (4) is possible. We have intentionally

given a derivation from first principles, based on the concept of a tangent as a line which intersects a curve in two coincident points, to show

that even if you do not remember a whole lot of formulas, things can

often be salvaged if you go by the basic principles.

A very straightforward problem about identifying the tangents to

a parabola passing through a given point. Perhaps the paper-setters

realised that the problem is too straightforward and hence given it a

twist by first making the candidates identify the points through which

the tangents are to pass. In the old days, instead of making the problem

numerical, it would have been asked to show that if two parabolas have

their vertices at the centre of an ellipse and their foci at the foci of that

ellipse then the tangents to either of them passing through the focus

of the other are mutually orthogonal. Such geometric results expressed

solely in words, have their own beauty.

Q.44 Let m and n be two

! positive integers greater than 1. If

n

ecos( ) e

e

m

lim

=

then

the

value

of

is

0

m

2

n

Answer and Comments: 2. As 0 both the numerator and

the denominator tend to 0. So the limit, say L, in the question if

0

of the indeterminate form . It is, therefore, tempting to apply the

0

LHopitals rule. But if m > 1 this rule will have to be applied again

and again. Let us temporarily assume that m = 1 and allow n to take

n

non-integral values as well. Define the function f () = ecos( ) 1.

Then LHopitals rule implies that

f ()

0

1

L = lim

(1)

Now, by a direct computation, for 6= 0 we have,

n

47

(2)

viz., sin(n ). This is valid for n > 0. But we have to be careful about

the last factor, viz. n1 . If n < 1, then this factor tends to as

n

0+ . The middle factor ecos( ) creates no problem as it tends to 1

when 0+ as long as n > 0. Nor does n which is a fixed number.

To see the result of the battle between the first and the last factor

when 0 < n < 1, let us rewrite (2) as

sin(n ) cos(n ) 2n1

f () =

e

(3)

depends on the last term, viz. 2n1 . If n > 1/2, then it will make

f () tend to 0. But if 0 < n < 1/2, then it will tend to as 0+ .

e

limit in question will exist. When it does, its value will be . In the

2

problem, of course, n is given to be an integer. But the argument above

m

has some fixed value, it is probably

suggests that if at all the ratio

n

2.

point. But a scrupulous person can make the argument valid as follows.

m

Let k = . Call m as . Then n = m/k and therefore, n = 1/k .

n

Moreover 0 as 0. So, by a change of variable,

L = lim

ecos(

1/k )

(4)

We are now back to the old game with replaced by and n replaced by

1/k. By the argument above, L will equal a finite non-zero value (which

will be necessarily e/2 ) only when 1/k = 1/2. This is equivalent to

saying that k = 2.

A mature person will approach the problem by considering the

comparable orders of magnitude. For example, when u 0, sin u is of

sin u

the same order as u because

tends to a finite non-zero limit as

u

cos u 1

tends

u 0. Similarly, cos u 1 is of the order of u2 because

u2

48

1

to the non-zero limit . Similarly, eu 1 is of the order of u. Now

2

factor out e from the numerator of the given expression and consider

n

ecos( )1 1

L = lim

(5)

0

m

Then clearly, L = eL provided L exists. By what we just said, the

numerator is of the order of cos(n ) 1, which, in turn, is of the order

of 2n . The denominator, on the other hand, is of the order of m . So

the ratio will tend to a finite non-zero limit only when the numerator

and the denominator have the same orders and that happens precisely

when m = 2n.

If need arises, this reasoning based on comparable orders of magnitude can be made precise by dividing and multiplying the expression

on the R.H.S. of (5) by cos(n ) 1 and then again by 2n . But that is

essentially a clerical work.

This is a very good, thought oriented problem. But since it is

a multiple choice question where no reason has to be given, a smart

candidate may get the correct answer by the sneaky path, i.e. by

assuming m = 1.

Q.45 If =

9x+3 tan1 x

12 + 9x2

x2 + 1

3

is

4

Answer and Comments: 9. This is plainly a question about evaluating a definite integral. Since the answer is to be an integer between

0 and 9, the last part operates as a hint to the value, viz. 1 + is of

the form ek+3/4 for some integer k. Although this does not help in

evaluating the integral, it serves to alert against numerical errors. In

the conventional form the question would have merely asked the value

of the integral.

The integral itself is easy if we use the substitution

u = 9x + 3 tan1 x

(1)

3

12 + 9x2

du

=9+

=

dx

1 + x2

x2 + 1

(2)

which yields

49

Hence we have

=

9+3/4

eu du

= e9+3/4 1

(3)

So ln(1 + ) = 9 3/4.

there is little choice about the correct substitution. The paper-setters

!

2

1

9 x 12 + 9x

at least could have given the integrand as (e )

(etan x )3 ,

x2 + 1

when in many other questions (for example, Q. 43 above) they have

given the data in an unnecessarily twisted form.

Q.46 Let f : IR IR be a continuous, odd function whichZ vanishes exx

1

f (t) dt for

actly at one point and f (1) = . Suppose F (x) =

1

Z x2

all x [1, 2] and G(x) =

t|f (f (t))| dt for all x [1, 2]. If

1

F (x)

1

lim

= , then the value of f (1/2) is

x1 G(x)

14

rule and the second Fundamental Theorem of Calculus (FTC). To apF (x)

ply the former to the limit lim

we must first ensure that it is

x1 G(x)

0

of the form. As both F and G are continuous (being functions de0

fined by integrals), this amounts to verifying that F (1) and G(1) both

vanish. This is not given. In fact, all we know is

F (1) =

and G(1) = =

1

Z 1

f (t) dt

(1)

t|f (f (t))| dt

(2)

(1) vanishes. For the same to happen for (2), we must ensure that

x|f (f (x))| is an odd function of x. For this we note that f (f (x)) is an

50

its absolute value |f (f (x))| is also an even function. Finally, x is an

odd function of x and so the product x|f (f (x))| is an odd function of

x, being the product of an even function and an odd function. So, the

integral on the R.H.S. of (2) vanishes as well.

We are now justified in applying the LHopitals rule. By the second

FTC, F (x) = f (x) and G (x) = x|f (f (x))| for all x [1/2, 2]. Hence

F (x)

F (1)

f (1)

F (x)

= lim

=

=

x1 G (x)

x0 G(x)

G (1)

1 |f (f (1))|

1/2

=

|f (f (1))|

1/2

=

|f (1/2)|

lim

1

. So, we now get

14

14

|f (1/2)| =

=7

2

(3)

(4)

We are still not quite through, because the problem asks for f (1/2)

and we only know that f (1/2) = 7. A sneaky reasoning would be

that since the correct answer has to be an integer from 0 to 9, 7 is

excluded. But, for an honest answer, we must show that f (1/2) > 0.

For this we apply the Intermediate Value Property (IVP). f (x) is given

to be continuous everywhere. Also f (1) = 1/2 > 0. So, if f (1/2) were

negative, then by the IVP, f would vanish at least once in the interval

(1/2, 1). But, being an odd function, f already vanishes at 0. As we

are given that f vanishes at exactly one point, we get a contradiction.

So, at long last we have fully justified our answer, viz. f (1/2) = 7.

To arrive at it legitimately, we had to give a lot of reasoning. First to

ensure the applicability of the LHopitals rule. Many students apply

this rule without ensuring this. Secondly, many students might not

bother to justify that f (1/2) > 0. And there is no way to tell if

they have really done the reasoning. So, this question rewards the

sloppy students and punishes the scrupulous ones. In the conventional

examinations, the two could be differentiated because reasonings would

have to be given.

51

rule and the FTC. Properties of odd functions have also been tested

along the way. But being an MCQ, all this fine creativity on the part

of the paper-setters is wasted. (In fact, the possibility that a student

might pick 7 as the correct answer simply because of the number 14

appearing in the problem cannot be entirely discounted. A safe gamble,

as there is no negative credit for a wrong answer.)

Q.47 Suppose that p~, ~q and ~r are three non-coplanar vectors in IR3 . Let the

components of a vector ~s along ~p, ~q and ~r be 4, 3 and 5, respectively.

If the components of this vector ~s along (~p + ~q + ~r), (~p q~ + ~r) and

(~p ~q + ~r) are x, y and z respectively, then 2x + y + z is

Answer and Comments: 9.. An extremely straightforward problem

about resolving a vector along three linearly independent vectors. We

are given

~s = 4~p + 3~q + 5~r

(1)

~s = x(~p + ~q + ~r) + y(~p ~q + ~r) + z(~p q~ + ~r)

= (x + y z)~p + (x y z)~q + (x + y + z)~r

(2)

(3)

As p~, ~q, ~r are non-coplanar and hence linearly independent, the coefficients in the two resolutions must match. This gives a system of three

equations in the three unknowns, viz.

x + y z = 4, x y z = 3, and x + y + z = 5

(4)

A straight substitution gives 2x + y + z = 8 + 1 = 9.

The word component in the statement of the problem is likely

to confuse some candidates. Sometimes, by a component of a vector ~s

along a vector ~p one understands, component along the direction of p~,

i.e. the component along a unit vector in the direction of p~. This usage

is common in expressions like the tangential and normal components

of acceleration used in physics. With this meaning, instead of (1) we

would have

4~p

3~q

5~r

~s =

+

+

(5)

||~p|| ||~q|| ||~r||

52

equation. This would make the problem not only more complicated,

but impossible to solve since the lengths of the vectors p~, ~q and ~r and

of their various linear combinations are not given.

But the paper-setters should not be blamed for this. The interpretation of component as they have intended is fairly standard. With the

alternate interpretation, the components are generally taken to be vectors along the respective directions. Nevertheless, the JEE authorities

have decided to award four points to every candidate for this question.

For those who interpret the word component as we have, the

question is a cakewalk. It is difficult to believe that such a simple

question is asked in an advanced test. It may be argued that this is

an atonement for some other difficult questions. But when these two

types of questions are combined, the selection takes place mostly on

the basis of such mediocre questions. That defeats the very purpose of

selection. And, for this, the paper-setters do deserve some criticism.

k

Q.48 For any integer k, let k = cos

7

Then the value of the expression

12

P

k=1

3

P

k=1

k

+ i sin

7

|k+1 k |

where i =

1.

is

|4k1 4k2 |

numbers lie at the vertices of a regular polygon with 14 sides inscribed

in the unit circle {z : |z| = 1}. (This happens because all the given

complex numbers are 14-th roots of unity.) For every k the vertices

k+1 and k are consecutive. So, |k+1 k | equals the length, say

d of each side of the heptagon. So the numerator of the expression is

12d. In the denominator also we are summing the distances between

consecutive vertices, except that in the denominator there are only 3

12d

terms. Hence the denominator is 3d. So the given ratio equals

= 4.

3d

A good problem based on the geometric representation of complex

numbers. Note that the number 7 has little role in the problem. It

could have been any positive integer r. Then the vertices would repeat

53

than 2r terms, the calculations would be the same.

The problem is reminiscent of a 1994 JEE problem which was a sort

of a converse to the present problem. There we were given a regular

n-gon with vertices A1 , A 2, . . . , An and the relationship between the

lengths of its sides and some of its diagonals. Specifically, it was given

that

1

1

1

=

+

A1 A2

A1 A3 A1 A4

(1)

a trigonometric equation. (Ironically, the answer to that problem was

also 7. See Comment No. 1 of Chapter 10.) However, the regularity

of the n-gon permits us to conclude from (1) a similar relationship for

any four consecutive vertices of the polygon, that is

1

1

1

=

+

Ai Ai+1

Ai Ai+2 Ai Ai+3

(2)

for every i, with the understanding that An+1 = A1 etc. This fact, along

with an application of what is called Ptolemys theorem about cyclic

quadrilaterals, gave an unexpected solution to the problem. (See again

Comment No. 2 of Chapter 10.) A solution using complex numbers is

also possible by taking the vertices to lie at the complex 2n-th roots of

unity. (See Comment No. 3 of the same chapter.) It is possible that

this solution has spurred the present problem. Of course, the present

problem is much simpler than the 1994 problem. How one wishes to

revive the golden days when such nice problems could be asked as full

length questions! All we have today is their weak imitations.

SECTION 2

This section contains EIGHT questions. Each questions has FOUR options (A), (B), (C) and (D). ONE OR MORE THAN ONE of these four

option(s) is(are) correct.

Marking scheme : +4 If fully correct, 0 if not attempted and 2 In all

other cases.

54

Q.49 Let f, g : [1, 2] IR be continuous functions which are twice differentiable in the interval (1, 2). Let the values of f and g at the points

1, 0 and 2 be given by the following table:

f(x)

g(x)

x = 1

3

0

x=0

6

1

x=2

0

-1

In each of the intervals (1, 0) and (0, 2) the function (f 3g) never

vanishes. Then the correct statement(s) is (are)

(A) f (x) 3g (x) = 0 has exactly three solutions in (1, 0) (0, 2)

(B) f (x) 3g (x) = 0 has exactly one solution in (1, 0)

(D) f (x) 3g (x) = 0 has exactly two solutions in (1, 0) and exactly

two solutions in (0, 2)

Answer and Comments: (B), (C). All the four statements are about

the zeros of h (x) where h(x) is defined by

(1)

(2)

Rolles theorem implies that h (x) vanishes at least once in each of the

intervals (0, 1) and (1, 2). If it had more zeros in either of these two

intervals, then by Rolles theorem again (applied to h (x) this time),

h (x) would vanish at least once. But we are given that this is not so.

Hence (B) and (C) are correct. That automatically precludes (A) and

(D).

A simple problem based on Rolles theorem. One fails to see what

was the idea of introducing the two functions f (x) and g(x) when the

entire problem is about the function h(x) = f (x) 3g(x). But then,

55

instead of the table, we would have to be given (2) and the paper-setters

probably thought that that would make the solution too obvious. So,

they worded the problem in terms of two functions. But nothing more

than elementary arithmetic is needed to convert the table to (2). The

question would have been interesting if some other pieces of information

were given from which (2) would follow, but not in such a direct manner.

(Contrast this with Q.46 above, where for an honest solution so many

properties of even and odd functions, and also the FTC had to be used.)

Since both the questions have equal marks, between two competing

candidates with severely limited time, the mediocre one who opts for

the present question will easily beat his possibly more intelligent and

scrupulous rival who might not have time left for this problem.

Q.50 Let f (x) = 7 tan8 x + 7 tan6 x 3 tan4 x 3 tan2 x for x (/2, /2).

Then the correct expression(s) is (are)

1

(A)

xf (x) dx =

12

Z 0/4

1

(C)

xf (x) dx =

6

0

Z

/4

(B)

(D)

/4

Z0

f (x) dx = 0

/4

f (x) dx = 1

Answer and Comments: (A) and (B). The paper-setters have carefully avoided the end points /2 from the domain because the tangent function is undefined at these points. Anyway that does not matter because in the problem the interval of integration is [0, /4] in all

statements. And there is no difficulty at the end points here. (If the

integrand tends to at either end points, the integral becomes what

is called an improper integral. Occasionally, improper integrals are

asked in the JEE. See Comment No. 16 of Chapter 18 and Exercise

(18.16). But they are not a part of the syllabus.)

Now, coming to the problem itself, the identity

tan2 x + 1 = sec2 x

(1)

f (x) = 7 tan6 x sec2 x 3 tan2 x sec2 x

56

(2)

The substitution

u = tan x

(3)

Z

/4

f (x)dx =

7u6 3u2 du

1

= u7 u3 = 1 1 = 0

0

(4)

The calculation of the integral in (A) and (C) is less direct. The same

substitution gives

Z

xf (x)dx =

(5)

As we already know an antiderivative of the second factor of the integrand, we resort to integration by parts and get

Z

/4

0

0

= 0

=

u2

1

(u7 u3 ) du (6)

+1

u3 (u2 1) du

u u5 du

1 1

1

=

=

4 6

12

(7)

problems asking a direct evaluation of some integral (instead of asking

some areas which first have to be written as integrals). The identity

(1), the substitution (3) and the factorisation of the last integrand in

(6) are so obvious, that a candidate who has cleared JEE Main can be

safely presumed to know them. So this problems tests no new ability.

It hardly belongs to an advanced test.

Q.51 Let f (x) =

If m

1/2

192x3

1

for all x IR with f ( ) = 0.

4

2

2 + sin x

f (x) dx M , then the possible values of m and M are

57

1

1

(B) m = , M =

4

2

(C) m = 11, M = 0 (D) m = 1, M = 12

(A) m = 13, M = 24

Answer and Comments: (D). We are not given f (x) directly. Instead we are given f (x) and f (1/2). In theory this determines f (x)

uniquely. In fact

f (x) =

1/2

f (t) dt =

1/2

192t3

dt

2 + sin4 t

(1)

the last integral in a closed form. Fortunately,

the problem does not

Z

ask for the exact value of the integral

1/2

estimates, i.e. some lower and upper bound on it. And these can be

obtained from some estimates on the integrand f (x), which, in turn,

can be obtained from properties of f (x) for x [1/2, 1]. Note that

f (x) is a ratio of two positive quantities. The numerator increases

strictly on the interval [1/2, 1] (in fact it increases on the entire real

line). But the denominator decreases on [1/2, 1] because of the term

sin4 x since the sine function decreases on the interval [/2, ]. So the

192

= 96 while its maximum

minimum of f (x) on [1/2, 1] is f (1) =

2+0

24

on [1/2, 1] is f (1/2) =

= 8.

2+1

Thus we have proved that

8 f (x) 96

(2)

for all x [1/2, 1]. From this, we now derive estimates for f (x) using the Fundamental Theorem of Calculus (FTC), and the fact that

f (1/2) = 0. We replace x by t in (2) and integrate over [1/2, x] to get

Z

1/2

8 dt

1/2

f (t) dt

1/2

96 dt

(3)

BY FTC, the middle integral is simply f (x) f (1/2) = f (x), while the

first and the last integrals are 8(x 1/2) and 96(x 1/2) respectively.

So we have

8x 4 f (x) 96x 48

58

(4)

Z

1/2

8x 4 dx

1/2

f (x) dx

1/2

96x 48 dx

(5)

1

1/2

1

48x)

= 1

= 24 12 = 12. These

1/2

are exactly the lower and upper bounds on the given integral as given

in (D).

This is a good problem which tests the ability to estimate when the

exact evaluation is not possible. Apart from some elementary knowledge of trigonometric functions, the entire solution is based on the fact

that inequalities are preserved by integrals, that is to say, if f (x) g(x)

Z

f (x) dx

for all x [a, b], then

a

a

of this kind and not those of the type of the last problem that are

suitable for an advanced test.

Q.52 Let S be the set of all non-zero real numbers such that the quadratic

equation x2 x+ = 0 has two distinct real roots x1 and x2 satisfying

the inequality |x1 x2 | < 1. Which of the following intervals is (are) a

subset(s) of S?

Answer and Comments: (A), (D). For the equation to be quadratic,

its leading coefficient

must be non-zero. When this is the case the

1 1 42

two roots are

. For the roots to be real and distinct, we

2

must have

42 < 1

Also the difference of the roots x1 and x2 is

1 42

|x1 x2 | =

||

(1)

(2)

1 42 < 2

59

(3)

(1/2, 1/2) while (3) holds if and only

only if

if (, 1/ 5)

(1/ 5, ). The intersection of these two

sets is (1/2, 1 5) (1 5, 1/2).

A very simple problem about the roots of a quadratic.

6

Q.53 If = 3 sin1 ( 11

) and = 3 cos1 ( 49 ), where the inverse trigonometric

functions take only the principal values, then the correct option(s) is

(are)

(A) cos > 0 (B) sin < 0 (C) cos( + ) > 0 (D) cos < 0

6

4

Answer and Comments: (B), (C), (D). Denote sin1 ( ) and cos1 ( )

11

9

by and respectively. Then = 3 and = 3. We are given

6

11

4

cos =

9

sin =

and

(1)

(2)

(1) and (2) using the formulas for the sines and cosines of triples of

angles. But the calculation of sin = sin 3 would require the value

4

of sin . It can be calculated from cos = . But that will involve a

9

radical and hence the calculation will be complicated. The calculations

for (C) will be even more complicated.

So we look for an easier alternative. If we look carefully, we are

interested only in the signs of the various expressions and not so much

in their exact values. And this can be done by locating the quadrants

in which the angles lie. We already know that both and lie in the

first quadrant, i.e.

0<<

2

2

(3)

3, 3 and 3( + ) lie. So we need sharper inequalities for and

than (3). How sharp should they be? If lies in an interval [a, b],

then 3 lies in [3a, 3b] whose length is thrice that of the first interval.

Since the angular measure of each quadrant is /2, to determine the

60

/6 or less.

4

Let us first tackle . We are given that cos = which is only

9

1

slightly less than = cos( ). As the cosine function is decreasing in

2

3

the first quadrant, we have

>

(4)

4

an upper bound, we need an angle whose cosine is slightly less than .

9

But here we need not be very choosy. Since we want to trap in an

interval of length /6, in view of (4), the upper bound /2 which we

already know from (3) will do. So we now have

<<

3

2

(5)

6

1

We are given that sin =

which is slightly bigger than . Since

11

2

1

sin(/6) = , we get /6 as a lower bound on . As for an upper

2

1

bound, we try /4 knowing that sin(/4) = . And, indeed, it works

2

6

1

because

does happen to be less than as we see by squaring both

11

2

the numbers. So we get

<<

6

4

(6)

We now have all the ammunition to fire the four shots. Since = 3,

because of (5) (, 3/2). So lies in the third quadrant where

both the sines and cosines are negative. This shows that (A) is false

while (B) is true. In a similar vein, because of (6), lies in (/2, 3/4)

which is a part of the second quadrant. So (D) is true. For (C), we

need both (5) and (6). Multiplying them by 3 and adding, we get

9

3

<+ <

2

4

61

(7)

The interval (3/2, 9/4) cuts across two quadrants, specifically, the

first and the fourth quadrant. But since the cosine function is positive

in both these quadrants, we see that (C) is true.

Like Q, 41 in Paper I, this question is also an excellent example of

testing the ability to quickly discard the most obvious approach. The

paper-setters have given just the right degree of hint by specifying that

all the inverse trigonometric functions are to be taken to have their

principal values.

Q.54 Let E1 and E2 be two ellipses whose centers are at the origin. The major

axes of E1 and E2 lie along the x-axis and the y-axis, respectively. Let

S be the circle x2 + (y 1)2 = 2. The straight line x + y = 3 touches

the curves S, E1and E2 at P, Q and R, respectively. Suppose that

2 2

PQ = PR =

. If e1 and e2 are the eccentricities of E1 and E2 ,

3

respectively, then the correct expression(s) is (are)

43

(A) + =

40

5

(C) |e21 e22 | =

8

e21

e22

7

(B) e1 e2 =

2 10

3

(D) e1 e2 =

4

in coordinate geometry. The equations of the ellipses E1 and E2 are

not given. We are only given that their major axes lie along the x- and

the y-axes respectively. Hence the equation of E1 is of the form

x2 y 2

+ 2 =1

a2

b

(1)

x2 y 2

+ 2 =1

c2

d

(2)

Let us first determine the points of contacts of the line (say L), x+y = 3

with S, E1 and E2 . Putting

y = 3x

62

(3)

which has x = 1 as its only root, confirming that L is a tangent to S.

From x = 1, y comes out to be 2 by (3). The point of contact, P is

P = (1, 2)

(4)

a, b, c, d.

Putting y = 3 x into (1) gives a quadratic

(a2 + b2 )x2 6a2 x + a2 (9 b2 ) = 0

(5)

a2 + b2 = 9

(6)

a2

3a2

=

. Hence the

Also the quadratic has only one root, viz. 2

a + b2

3

point of contact comes out to be

a2 9 a2

,

)

3

3

Q=(

(7)

2 2

As we are given P Q =

, we get

3

(

9 a2

8

a2

1)2 + (

2)2 =

3

3

9

(8)

from (6), b2 = 4 or 8. But we are given that a > b. So we have

a2 = 5, b2 = 4. Since the eccentricity e1 of the first ellipse is given by

b2 = a2 (1 e21 )

(9)

we get

e1 =

1

4

=

5

5

(10)

respectively. Since P Q = P R, (8) also holds with a replaced by c. So

63

time we are assuming d > c. So we have to take the second possibility,

viz. c = 1 and d = 8. Then e2 is given by

s

s

c2

7

1

e2 = 1 2 = 1 =

(11)

d

8

2 2

Now that we know both e1 and e2 , by a direct calculation we get

7

e1 e2 =

(12)

2 10

43

1 7

+ =

(13)

e21 + e22 =

5 8

40

1 7

27

and finally, |e21 e22 | = | | =

(14)

5 8

40

The calculations in the problem are straightforward. Also the

symmetry of the data allows us to apply most of the work for the

first ellipse to the second ellipse. The only difference comes in the

last step because of the interchange of the directions of the major and

the minor axes. The problem would have been a good one as a full

length question in a conventional test. In the present set-up, the time

allotted is far too inadequate. In this problem, there does not seem

to be any way to directly compute, say e1 e2 , without first identifying

e1 and e2 individually (as there would have been if we could write a

quadratic whose roots are e1 and e2 ). So, logically, the problem ends

with the determination of e1 and e2 . The further calculations needed

in finding the numerical values of the various expressions in the four

options are pure drudgery and test no valuable quality on the part of

the candidates.

Q.55 Consider the hyperbola H : x2 y 2 = 1 and a circle S with centre

N(x2 , 0). Suppose H and S touch each other at a point P (x1 , y1 ) with

x1 > 0, y1 > 0. The common tangent to H and S at P intersects the

x-axis at a point M. If (l, m) is the centroid of the triangle P MN,

then the correct expression(s) is (are)

64

1

dl

= 1 2 for x1 > 1

dx1

3x1

1

dl

= 1 + 2 for x1 > 1

(C)

dx1

3x1

dm

x1

= q

for x1 > 1

dx1

3 x21 1

dm

1

(D)

= for y1 > 0

dy1

3

(B)

(A)

Answer and Comments: (A), (B), (D). The notations in the problem

are slightly confusing. Usually, P (x1 , y1) would denote a fixed point in

the plane. In the present problem, it is a variable point in the first

quadrant which moves on the hyperbola H. But instead of beginning

the problem by saying so, the point (x1 , y1 ) is specified as the point

of contact of this hyperbola with a circle S whose centre is given as

N = (x2 , 0). So, it is to be assumed that x2 is also not fixed. So, the

picture we gather is that x1 and y1 are functions of the variable x2 .

But the problem asks for derivatives w.r.t. the variables x1 and y1 .

y

H

S

x

O

express N in terms of P . In fact, the circle S has little role to play

except to tell us that the line P N is normal to S and hence also to H

at P , because S and H have a common tangent at P . So, we could

as well define N as the point where the normal to H at P meets the

x-axis. The point M, on the other hand, is specified directly in terms

of (x1 , y1), as the point where the tangent to H at P cuts the x-axis.

So the essence of the data of the problem is this. Suppose P (x1 , y1 )

is a variable point in the first quadrant portion of the hyperbola H :

x2 y 2 = 1, M and N are, respectively, the points where the tangent

65

and the normal to H at P cut the x-axis. And finally, (l, m) is the

centroid of the triangle P MN. The coordinates of M and N and

hence those of the centroid will be functions of the variables x1 and y1 .

But these two variables are not independent of each other. Either of

them can be expressed in terms of the other by the equations

y1 =

and x1 =

x21 1

(1)

y12 + 1

(2)

The same goes for m. The problem then asks about the derivatives of

these functions.

Let us now turn to the problem itself. The equation of the tangent

to H at (x1 , y1 ) is

xx1 yy1 = 1

Its point of intersection with the x-axis is (

M =(

(3)

1

, 0). So, we have

x1

1

, 0)

x1

(4)

x1

y1

, that of the normal at P is .

y1

x1

Hence the equation of the normal to H at P is

As the slope of the tangent is

y y1 +

y1

(x x1 ) = 0

x1

(5)

N = (2x1 , 0)

(6)

1

1

1

(3x1 + ) = x1 +

3

x1

3x1

y1

and m =

3

l =

66

(7)

(8)

1

dl

= 1 2 which shows that (A) is true and (C)

dx1

3x1

is false. Also from (8), we see that (D) is true. For (B), we first need

to express m as a function of x1 . For this we combine (8) with (1) to

get

From (7) we get

m=

x21 1

3

(9)

Differentiating,

x1

dm

= q

dx1

3 x21 1

(10)

wonders if it was intentional. If so, the idea was probably to test a

candidates ability to grab the essence of a problem from a nebulous

data. But it does look silly. It is like a person telling you that he met

his fathers wifes daughter, instead of saying that he met his sister

(possibly a step-sister)!

Q.56 The option(s) with the values of and L that satisfy the following

equation is (are)

Z

Z0

et (sin6 t + cos4 t) dt

t

=L

e (sin t + cos t) dt

e4 1

(A) = 2, L =

e 1

e4 1

(C) = 4, L =

e 1

e4 + 1

(B) = 2, L =

e +1

e4 + 1

(D) = 4, L =

e +1

two integrals, say I and J. Both the integrals have the same integrand,

say

f (t) = et (sin6 t + cos4 t)

67

(1)

I = I1 + I2 + I3 + I4

(2)

where

I1 =

f (t)dt, I2 =

f (t)dt I3 =

f (t)dt and I4 =

4

3

f (t)dt (3)

the periodicity of the trigonometric functions. The functions sin2 t

and cos2 t and hence all their powers are periodic with a period of

for both = 2 and = 4 (in fact, for all integral values of ) because

sin((t + )) = (1) sin t

and cos((t + )) = (1) cos t

(4)

(5)

for all t. However, the first factor of the integrand, viz. et is not

periodic. Still we have the identity

et+ = e et

(6)

f (t + ) = e f (t)

(7)

I2 =

=

=

Z

Z0

0

f (t) dt

f (u + ) du

e f (u) du

= e I1

(8)

R.H.S. of (2) are in a geometric progression with common ratio e and

the first term I1 . Hence

I = I1

e4 1

e 1

68

(9)

This holds regardless of what the integer is. So both (A) and (C) are

correct.

Problems involving integrals of periodic functions have been asked

in the JEE (e.g. see Exercise (18.8)). But the present one is a novel

one in that it combines periodic functions and a simple property of the

exponential functions. Once the idea strikes, the computations involved

are minimal. So this is a very good problem.

Although it is well beyond the JEE level, we remark that similar

calculations arise in finding what is called the Laplace transform of

a periodic function. Suppose f (t) is a continuous function defined for

all t 0. Then assuming that the growth of |f (t)| as t is not too

wild, it can be shown that the improper integral

F (s) =

est f (t)dt

(10)

We now have a new function, F (s), of a new variable s. It depends on

the function f (t) and is often denoted by L(f ). As easy examples, we

invite you to show, using integration by parts, that

L(tn ) =

n!

sn+1

, n = 0, 1, 2, . . .

1

s

s

L(cos t) = 2

( IR)

s + 2

( IR)

and L(sin t) = 2

s + 2

L(et ) =

(11)

(12)

(13)

(14)

What makes the Laplace transforms very useful is the result, easily

proved by integration by parts, that for a continuously differentiable

function f (t),

L(f ) = sL(f ) f (0)

(15)

As a consequence, a differential equation for a function f (t) is transformed into an algebraic equation for its transform F (s). Algebraic

equations are easier to solve. So just as logarithms are useful in arithmetic because they transform the multiplication of two positive real

69

numbers (which is time consuming) into the addition of their logarithms (which is much easier), Laplace transforms are very useful in

solving differential equations.

Now suppose f (t) is a periodic function with a period T > 0. That

means

f (t + T ) = f (t)

(16)

for all t IR. Then by taking steps similar Zto those in the solution of

the present problem, for every s, the integral

0

for n = 1, 2, 3, . . .. As a result, we get

F (s) = I1 + I2 + . . . + In + . . .

(17)

where

In =

nT

(n1)T

est f (t)dt

(18)

(17) and (18) are valid for any f (t). But if f (t) is periodic with a period

T , then by taking steps similar to the derivation of (8) above, we get,

that for every fixed s,

In+1 =

(n+1)T

nT

Z nT

(n1)T

nT

(n1)T

sT

e In

est f (t) dt

es(u+T ) f (u + T ) du

esT esu f (u) du

(19)

esT . So the the R.H.S. of (17) is an infinite geometric series. As a

result, we get

F (s) =

I1

1 esT

70

(20)

the integral I1 is not improper. So periodicity of f (t) allows us to

express its Laplace transform L(f ) in terms of a proper integral over

an interval whose length equals the period of f (t).

The JEE paper-setters are undoubtedly familiar with Laplace

transforms and chances are that the present problem was spurred as a

miniature version of (20).

SECTION 3

This section contains TWO paragraph. Based on each paragraph, there

will be TWO questions Each question has FOUR options (A), (B), (C) and

(D). ONE OR MORE THAN ONE of these four option(s) is(are) correct For

each question, darken the bubble(s) corresponding to all the correct option(s)

in the ORS.

Marking scheme :

+4 If only the bubble(s) corresponding to all the correct option(s) is(are)

darkened

0 If none of the bubbles is darkened

-2 In all other cases.

PARAGRAPH 1

Let F : IR IR be a thrice differentiable function.

Suppose that F (1) = 0, F (3) = 4 and F (x) < 0

for all x (1/2, 3). Let f (x) = xF (x) for all x IR.

(A) f (1) < 0

(C) f (x) 6= 0 for any x (1, 3)

(D) f (x) = 0 for some x (1, 3)

Answer and Comments: (A), (B), (C). The problem asks about

f (x) and its derivative f (x). We are given

f (x) = xF (x)

71

(1)

f (x) = F (x) + xF (x)

(2)

(3)

Hence

since F (x) < 0 for all x (1/2, 3). So, (A) is true.

For (B),

(4)

We are not given the values of either F (2) or F (2). But we are given

that F (x) < 0 for all x (1/2, 3). So F (x) is strictly decreasing

on [1/2, 3]. (This assertion requires the use of Lagranges Mean Value

Theorem. But it is possible that this subtlety is missed by many candidates, and, being an MCQ, there is no way to tell.) Since F (1) = 0

and 2 > 1, we get F (2) < F (1) = 0. Also since 2 (1/2, 3), F (2) < 0

straight from the hypothesis. So by (3), (B) is also true.

The options (C) and (D) are the logical negations of each other.

So exactly one of them is true. The reasoning given in (B) for x = 2

applies for any x (1, 3) and so from (2), f (x) < 0 for all x (1, 3),

as both the terms on the R.H.S. are negative. So (C) is true, and

automatically, (D) is false.

Q.58 If

x2 F (x) dx = 12 and

pression(s) is (are)

(B)

(D)

Z1 3

1

f (x) dx = 12

f (x) dx = 12

Answer and Comments: (C), (D). All options involve exact values of some integrals and not just their sign determinations. So, the

inequalities given in the paragraph are not likely to be of much help.

From the nature of the integrands, it is obvious that integration by

72

f (x)dx =

xF (x) dx nothing

an antiderivative of the first factor, viz. x of the integrand. We then

get

Z

f (x) dx =

xF (x) dx

Z 3 2

3

x2

x

= ( F (x))

F (x) dx

1

2

1 2

(5)

9

1

substitution, the first term is F (3) F (1). The values of F (3) and

2

2

F (1) are given as 4 and 0 respectively. Substituting these, the R.H.S.

of (5) comes out to be 18 0 + 6 = 12. Hence (D) is correct and

(B) false.

To determine which (if any) of the remaining options is true, we use

the second piece of information, viz.

Z

x3 F (x) dx = 40

(6)

Z

1

3x2 F (x) dx

(7)

Again, we are given the value of the integral on the R.H.S. as 36. So,

by a direct substitution we get

27F (3) F (1) = 40 36 = 4

(8)

Since the options (A) and (C) are in terms of f rather than F , we

must convert F (3) and F (1) suitably. Applying (2) and the fact that

F (3) = 4, we have

F (3) =

f (3) F (3)

f (3) + 4

=

3

3

(9)

f (1) F (1)

F (1) =

= f (1)

1

73

(10)

9f (3) + 36 f (1) = 4

(11)

A simple, but extremely weird pair of questions. Conceptually, the

two questions have little in common. But because of the constraint in

paper-setting, they have been artificially clubbed together into a single paragraph. The first question is based on applying the Lagranges

Mean Value Theorem to get some inequalities. The second one is based

on integration by parts. The work involved is repetitious and the arithmetic required is arbitrary and prone to numerical errors.

PARAGRAPH 2

Let n1 and n2 be the number of red and black balls,

respectively, in box I. Let n3 and n4 be the number

of red and black balls, respectively, in box II.

Q.59 One of the two boxes, box I and box II was selected at random and a

ball was drawn at random from out of this box. The ball was found to

be red. If the probability that this red ball was drawn from box II is

1

, then the correct option(s) with the possible values of n1 , n2 , n3 and

3

n4 is (are)

(A) n1 = 3, n2 = 3, n3 = 5, n4 = 15

(B) n1 = 3, n2 = 6, n3 = 10, n4 = 50

(C) n1 = 8, n2 = 6, n3 = 5, n4 = 20

(D) n1 = 6, n2 = 12, n3 = 5, n4 = 20

Answer and Comments: (A), (B). This is a problem of conditional

probability. We are first selecting a box and then drawing a ball from

74

the selected box. We want to find the probability that the ball drawn

is from the second box, given that its colour is red.

There is a chain of two processes here, first selection of a box and

then the draw of a ball. So for a visual representation of the thought

process, tree diagrams are more convenient than Venn diagrams as

elaborated in Comment No. 11 of Chapter 22. The appropriate tree

diagram for the present problem is shown below.

A

1/ 2

1/ 2

II

I

p1

p

4

Here A is the starting node. From there, there are two branches leading to the two nodes I and II corresponding to the two boxes. Each

branch has probability 1/2. From each of these two nodes, there are

two branches leading to the nodes marked r and b depending upon the

colour of the ball drawn. Their probabilities are marked as p1 , p2 , p3 , p4 .

They can be easily calculated from the data of the problem as

p1 =

n1

n2

n3

n4

, p2 =

, p3 =

and p4 =

n1 + n2

n1 + n2

n3 + n4

n3 + n4

(1)

Now let R be the event that the ball drawn is red. Then from the

diagram above,

1

1

p1 + p3

2

2

1 n1

1 n3

=

+

2 n1 + n2 2 n3 + n4

P (R) =

(2)

Let E2 be the event that the second box was chosen. Of course P (E2 ) =

1

. But the desired probability is not this, but rather the conditional

2

75

probability P (E2 |R), i.e. the probability that the ball drawn is from

the second box given that it is red. By Bayes theorem (see Comment

No. 5 of Chapter 22), or more easily, by common sense,

P (E2 |R) =

P (E2 R)

P (R)

(3)

numerator is simply the second term of the R.H.S. of (2). So, the

desired probability, say p, is

n3

n3 + n4

(4)

p=

n3

n1

+

n1 + n2 n3 + n4

The conceptual part of the problem ends here. In a numerical

problem, we would be given the values of n1 , n2 , n3 and n4 and asked

1

to calculate p from (4). But here we are given that p =

and are

3

asked to see for which of the given sets of values of n1 , n2 , n3 and n4 ,

1

the answer comes out to be . So, in essence, we are asked to repeat

3

the arithmetical work four times. A slight simplification is possible.

p3

1

Note that p can be written as

. So this will equal if and only

p1 + p3

3

if p1 = 2p3 . In (A), p1 = 1/2, p3 = 1/4 and so (A) is true. In (B),

p1 = 1/3, p2 = 1/6 and so (B) is also true. In (C), p1 = 4/7, p3 = 1/5

and so (C) is false. Finally, in (D), p1 = 1/3, p3 = 1/5 and so (D) is

false.

A very standard problem of computing conditional probability.

But one wonders what is gained by forcing the candidate to do the

last bit of arithmetic four times. Why not give only one set of values

of n1 , n2 , n3 and n4 and give four numerical alternatives for the value

of p? This used to be the standard practice in the past. But there

have been cases where, because of a carelessness on the part of the

paper-setters, the fake answers could be eliminated easily without doing

the problem honestly, for example, when the fake answer was given as

a rational number whose denominator did not divide the size of the

sample space. It is probably to avoid this kind of a criticism that this

time the paper-setters have chosen to torture the candidates.

76

Q.60 A ball is drawn at random from box I and transferred to box II. If the

1

probability of drawing a red ball from box I, after this transfer, is ,

3

then the correct option(s) with the possible values of n1 and n2 is(are)

(A) n1 = 4 and n2 = 6

(B) n1 = 2 and n2 = 3

(C) n1 = 10 and n2 = 20 (D) n1 = 3 and n2 = 6

Answer and Comments:. (C), (D). It is given that a ball, randomly

drawn from box I, is transferred to box II. But the second part of the

problem deals only with what is left in box I. So what happens to

this ball after its removal from box I is irrelevant. The paper-setters

could simply have said that one ball at random was removed from box

I, or kept the data as it is except that the second ball is drawn from

box II. That would make the transfer relevant and the problem more

interesting.

Let p1 and p2 be the probabilities that the ball removed was red

and black respectively. Clearly,

p1 =

n1

n2

and p2 =

n1 + n2

n1 + n2

(5)

Now, if the ball removed is red, then box I is left with n1 1 red and

n2 black balls. The probability of drawing a red ball from this box is

p3 , given by

p3 =

n1 1

n1 + n2 1

(6)

Similarly, if the ball removed is black, then box I is left with n1 red

and n2 1 black balls and the probability, say p4 , of drawing a red ball

from it now is

n1

p4 =

(7)

n1 + n2 1

With the help of a tree diagram, if necessary, the probability, say p,

that the second ball drawn from box I is red, is

p = p1 p3 + p2 p4

n1 (n1 1)

n2 n1

=

+

(n1 + n2 )(n1 + n2 1) (n1 + n2 )(n1 + n2 1)

77

n1 (n1 + n2 1)

(n1 + n2 )(n1 + n2 1)

n1

=

n1 + n2

(8)

1

. This will happen if and only if n2 = 2n1 .

3

Clearly this is the case in options (C) and (D) but not in the others.

We are given that p =

This problem is much simpler than the last one. In fact, a perceptive

student will hardly fail to notice that p is the same as p1 , i.e. the same

as the probability of drawing a red ball from the box I at the start. It is

unlikely that he will have the time to ponder if this is just a coincidence

or can be somehow justified. But if he does, he can paraphrase the

problem to say that two balls were drawn from box I, one after the

other, and the problem asks for the probability that the second ball is

red. How do we justify that this is the same as the probability that the

first ball is red?

Intuitively, the answer is that we can interchange the two balls

drawn. For a formal proof, let n = n1 + n2 . Let us suppose that the

red balls as well as the black balls are distinct from each other. (This

does not affect the probabilities of the events involved.) Let X be the

set of all these n balls. Our sample space, say S is the set of all ordered

pairs (x, y) X X such that x 6= y. Clearly S has n(n 1) elements.

Now let F1 and F2 , respectively, be the favourable subsets for the event

that the first ball is red and the event that the second ball is red. In

symbols,

F1 = {(x, y) S : x is red}

and F2 = {(x, y) S : y is red}

(9)

(10)

As the sample space is the same, proving that the two events have the

same probability is equivalent to showing that the sets F1 and F2 have

the same number of elements. This can be done by defining a bijection

f : F1 F2 which interchanges the two balls drawn. Formally,

f (x, y) = (y, x). Then f is its own inverse and hence a bijection.

Note that this argument will also apply if we draw 3 balls, one

after the other without replacement. The probability that the first one

be red equals the probability that the second one be red and also the

78

probability that the third one be red. Indeed this will hold even if the

box contains balls of more than two colours and we draw any number

of balls from it as long as this number does not exceed the total number

of balls in the box. (If it does, then the sample space is vacuous and

the probability is undefined.)

As in every year, the MCQ format makes it impossible to reward a

candidate who has the ability to do this extra bit of thinking. In fact,

he is advised against it because it is a sheer waste of his precious time.

79

CONCLUDING REMARKS

The paper-setters have been careful to avoid mathematical mistakes and

ambiguities. For example, in both the questions involving the composite

of two functions, the order in which the two functions act has been clearly

specified. In Q. 47 of Paper 1, the symbol a is used to denote a variable.

This should have been made clear, because normally the symbols a, b, c etc.

are used for constants and x, y, z, t etc. for variables. Similarly, in Q. 55 of

Paper 2, it should have been made clear that x2 and x1 are variables. In

Q.54 of Paper 1, instead of saying true it would have been better to say

necessarily true to avoid the degenerate possibility of a zero matrix which

is symmetric as well as skew-symmetric. But these omissions are unlikely to

have caused any serious confusion.

Except for a slight lapse of the article in Q. 58 of Paper 2, (saying, a circle

with a given diameter instead of saying the circle) there are no anomalies

in the wording of the problems. As pointed out in Q.60 of Paper 2, the information that the ball removed from box I was put into box II is irrelevant.

But again, these are things to be forgotten after a slight chuckle. On the

contrary, as commented in Q. 47 of Paper 1, the paper-setters have commendably avoided an inconsistency in the data which could have resulted if

there was a repetition of carelessness in the drafting of a similar question in

JEE 2011.

A controversy did arise about the word component used in Q.47 of Paper

2. In a conventional examination, from the answerbook of a candidate, it

would be possible to see which interpretation he had taken. And then a

decision could be taken as to how much credit be given. The MCQ format

makes this impossible. The JEE organisers have decided to award full marks

to every candidate. As a result, those who interpreted and solved the problem

correctly are at a loss. Fortunately, it is not a serious loss because the

problem was very trivial anyway. The JEE Advanced organisers deserve to

be commended for showing this openness of mind. (In an earlier version of

this commentary, the JEE Advanced organisers were wrongly criticised. The

mistake is regretted.)

Given the constraints imposed on them, the paper-setters have come up

with some good problems. The two that stand out are Q.41 in Paper 1

(about the reflection of a parabola) and Q.53 of Paper 2 (about the sign

determination of the trigonometric functions of angles that are triples of some

given angles). Q. 51 of Paper 1 (about the differentiability of the composite of

80

two functions), Q.44 (about orders of magnitude) and Q.46 (about functions

defined by integrating odd functions) of Paper 2 are also very good questions

of the conceptual type. But their ability to separate the men from the boys

is marred by the multiple choice format.

Some other good questions in Paper 1 are Q.43 (about boys and girls

standing in a row), Q.47 (about a function defined by integrating cos2 t), Q.54

(about symmetric and skew-symmetric matrices), Q.55 (about expanding a

determinant of order 3) and Q.59(C) (about the complex cube roots of unity),

while those in Paper 2 include Q.41 (about the sums of the terms in an

arithmetic progression), Q.48 (about representing complex roots of unity as

vertices of a regular polygon), Q.51 (about finding lower and upper bounds on

an integral), Q.54 (about the eccentricities of two ellipses) and Q.56 (about

integrating products of periodic functions with the exponential function).

Questions which are too straightforward or familiar and hence have no

place in an advanced test include Q.42 (about complementary probability

with a coin toss), Q.46 (about minimising the material needed for a cylindrical container), Q.49 and Q.50 (both on differential equations) and Q.56

(about parametrised family of planes) in Paper 1 and Q.47 (about resolving a vector), Q.49 (about applications of Intermediate Value Property and

Rolles theorem), Q.50 (a straight, simple integral), Q.52 (about the roots of

a quadratic) and both the probability problems (Q.59 and 60) in Paper 2.

There are also numerous questions which reduce to writing down and solving

a system of equations, often in two unknowns. As there is nothing very great

in writing these equations or in solving them, such questions only reward the

mediocrity.

Many of these questions could have been dropped and room made to

accommodate certain areas that are singularly absent, such as number theory,

surds, logarithms and inequalities. This is probably the first year when

there is no question based on the A.M.-G.M. inequality. It is true that

the multiple choice format precludes all proofs and that includes proofs by

induction, combinatorial identities, triangular optimisation and many results

in theoretical calculus. Still, in the past the paper-setters have managed to

give at least a token place to the areas just mentioned.

There is also a tendency apparent in many questions to force the candidates to do extra, repetitious drudgery after the conceptual part of the

problem is over. For example, Q.54 in Paper 2 logically ends with the determination of the eccentricities of the two ellipses. Subsequent work is sheerly

clerical. Similarly, in Q.59 and 60 about probabilities, instead of giving some

81

particular values of the parameters n1 , n2 , n3 and n4 and asking for a numerical answer, which is the standard practice, the paper-setters have given four

sets of such values and asked the candidates to identify those sets of values

in which the answer will come out to be a given number. One really fails

to see the purpose of such a twist when what it tests is not even remotely

related to any probability, It is sheer arithmetical drudgery. In Q.53 of Paper

1, once the vectors ~a, ~b and ~c are identified correctly, little purpose is served

by asking the candidates to compute the various expressions involving their

dot or the cross products.

Some questions demand work that is far in excess of what is justified by

their credit. A notable example is Q.60(D) where a candidate has to sketch

three regions and find their areas all for 2 marks.

There is no significant difference either in the standard or the topics covered in the two papers. One really wonders the rationale behind having two

papers each having all three subjects. Why not have three separate papers,

one for each subject and give the paper-setters of each subject the freedom

to choose the format of that paper as demanded by the peculiarities of the

subject, rather than impose the same format for all the three subjects? In

Chemistry, many of the questions are memory oriented. You simply have to

know the answer beforehand. There is no way you can get it by deductive

reasoning. For example, questions that ask for the compounds that result

from some chemical reactions.

As a sample, as many as eight questions out of the 20 questions in Chemistry in Paper 2 of this years advanced JEE directly ask to identify the

products of some reactions. Q.39 in Paper 1 is for eight marks and asks to

identify the radicals (from a given list) occurring in each of the five given

ores, viz. siderite, malachite, bauxite, calamine and argentite. This is sheer

memory. If you know the chemical composition of these ores you get 8 marks

within a few seconds. If you dont, just forget the question. No amount

of deductive reasoning can help you. Its counterpart in mathematics, viz.,

Q.59 in Paper 1, requires you to solve four completely unrelated problems,

each requiring some thought and moderate computation. There is simply no

comparison between the two.

One can argue that even in chemistry, there are questions where some

thought and numerical work is needed. This is true. But the thrust in these

questions is on knowing some reaction. For example, in one of the questions

in Paper 2, a closed vessel with rigid walls containing 1 mol radioactive

uranium and one mol of air is given. The question then asks what will be

82

the ratio of the final pressure to that of the initial one after the uranium

decays completely. The answer comes from the fact that in the radioactive

decay, every atom of uranium produces one atom of lead and eight atoms of

helium. This is pure memory. The thought needed afterwards is that now

in the vessel we have 8 mols of gas in addition to the original 1 mol of air.

As the volume and the temperature are the same, the new pressure will be

9 times the original. This bit of mathematics is childishly simple. The same

thing is true about a question in which diborane reacts with methanol.

The mathematical equivalent of such a single idea question, would be, for

example, to give a right angled triangle with hypotenuse 5 and one side 4 and

ask to find the third side, or to ask in how many ways

5 boys and5 girls can

stand

in

a

row.

In

both

the

cases,

the

answers

(viz.

52 42 = 25 16 =

But such questions are not asked in mathematics even in the elimination

round. The second question appears as a part of Q.43 in Paper 1 this year.

But it is only the beginning and not the end. In fact, the real problem then is

to count in how many of these 10! arrangements, all the 5 girls stand together

and also in how many of them exactly four girls stand together. This cannot

be done just by a hand counting. It requires a thought. Once the correct

thought is grasped, the arithmetic needed is elementary. But the thought

has to come from the candidate. No manual can help.

This is certainly not to deride Chemistry as a discipline of science. But

the point is to decide which qualities are more important when they are being

tested by an examination like the JEE whose avowed purpose is to select the

most gifted candidates. It would be great if the IITs get a sufficient number

of candidates who are good in all respects. That is, they know the chemical

composition of bauxite as well as have the ability to solve the problem above

of 5 boys and 5 girls. Unfortunately, this does not happen and so the IITs

have to make a choice. So, suppose the IITs have to choose between two

students, say A, who can do the problem of 5 boys and 5 girls but does not

know, or has forgotten what bauxite is, and another candidate B who knows

what bauxite is but cannot do the thinking needed for the problem of 5 boys

and 5 girls. Whom should they choose? Clearly, the choice should be A,

because he can be taught what bauxite is, or, when need arises, can easily

look it up. But, although B can also be taught how to do this particular

problem, the ability to do such thinking originally cannot be inculcated into

him as it is largely an innate quality.

In Physics the picture is better than in Chemistry. There are virtually no

83

questions which are purely memory oriented. Most questions require some

work to arrive at the answer. But usually, this work consists of applying one,

or in some cases, two standard formulas. Once you know the right formulas

to apply, the rest of the work is purely mathematical. For example, in one

of the Physics questions in Paper 2 this year, two solid spheres of equal

radii R but with different mass density functions are given and we have to

find the ratio of their moments of inertia about axes passing through their

respective centres. This is a standard application of integrals. At one time it

was taught in mathematics under mechanics. There is one question where a

parallel plate capacitor is given, the plates being separated by two dielectric

slabs of equal thickness but different relative permittivities. We are asked

to find its capacitance given its capacitance in air. A diagram is also given.

That makes it abundantly clear that we have to apply the formula for the

capacitance of two capacitors in series. Once this is realised, the mathematics

involved is straightforward.

In fact, in some questions in the Physics paper, the relevance with physics

is superficial. Take, for example, this question from Paper 2 where we are

given that the energy of a system as a function of time t is E(t) = A2 et ,

being a given constant. We are further given the percentage errors in the

measurements of A and of t and are asked to calculate the percentage error

in E(t) at t = 5. Now, this is a purely mathematical question, involving

first order approximation using derivatives. (Ironically, that is usually done

in the first mathematics course in the IITs after the successful candidates

enter them!) There is hardly any physics in this problem except the name.

The problem could as well have been posed as a problem in economics by

replacing the phrase the energy of a system at time t by a phrase like the

demand for some commodity at time t.

The point to note is that once a question in chemistry is understood, it

takes very little time to answer it, provided you know the answer. In physics,

it takes some work to arrive at the answer. But that work is usually some

standard mathematical computation. Moreover, the physics paper-setters

usually make it easier for the candidates to understand the question, by

supplying appropriate diagrams. This year, for example, the Physics section

of Paper 2 has as many as 10 diagrams for the 20 questions.

Contrast this with mathematics. There is not a single diagram in the

mathematics sections of either of the two papers! And this year is not an

exception. In fact, some questions seem to be intentionally obscure. For

example, Q.55 of Paper 2 about the tangent and the normal to a hyperbola.

84

This problem as well as many other problems where the data is geometric

would be easier to understand if appropriate diagrams were given. But then

the trouble is that they will also be easier to solve. For example, if a diagram

were given for Q.41 of Paper 1 (about the reflection of a parabola) that would

also give an unwarranted hint to the solution. And that will kill the beauty

of the problem. One of the qualities considered desirable in mathematics is

the ability to read between the lines.

Summing up, Chemistry questions test the knowledge of facts, Physics

questions test the ability to apply some standard formulas. But Mathematics

is inherently different. A mathematical equivalent of the question about the

capacitor considered above, would be, for example, a question asking for

the area of a triangle, given its two sides and one angle. None of the good

questions listed above are of this type. There is an element of art and sport

in mathematics which is generally absent in physics and chemistry which are

sterile sciences. It becomes very difficult for the paper-setters in mathematics

to cater to this element if they are forced to frame the questions in the

same format as physics and chemistry. Designing and solving a truly good

mathematics problem is an artistic experience. No wonder some of the good

problems in the old JEEs are remembered even after decades. (Two such

problems from the 1994 JEE have been referred to in this commentary, one

in Q.55 of Paper 1 and the other in Q.48 of Paper 2.) Can that be said about

any of the mathematics problems of the recent JEEs?

True improvement would come only if the unique position of mathematics

is recognised by the policy makers. When the candidates are already screened

by JEE Main and the number of those eligible for the Advanced test is reduced to a fairly manageable size, there is no reason to duplicate the selection

through another test of the same level and the same format. The Advanced

JEE ought to be qualitatively different from the Main one. At present it is

not. Because of the constraints on the paper-setters in mathematics, they

are forced to ask questions which fit only in the screening round but hardly

qualify for the final selection.

Ideally, at the Advanced JEE, there should be two papers of three hours

each, one solely for mathematics and the other for a combination of physics

and chemistry in the proportion, say 2 : 1. Since the Main examination

already tests the degree of exposure to the three subjects, there should be no

compulsion to cover the entire syllabus in the Advanced papers. Moreover,

mathematics paper-setters should have the freedom, as they did many years

ago, to design questions requiring different amounts of work and allot the

85

of thinking of a candidate can be assessed, as he has to give reasoning for his

work.

One can only pray that this will happen in near future.

86

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