Académique Documents
Professionnel Documents
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Economics
Andrew McLennan
April 8, 2014
Preface
Over two decades ago now I wrote a rather long survey of the mathematical
theory of fixed points entitled Selected Topics in the Theory of Fixed Points. It
had no content that could not be found elsewhere in the mathematical literature,
but nonetheless some economists found it useful. Almost as long ago, I began
work on the project of turning it into a proper book, and finally that project is
coming to fruition. Various events over the years have reinforced my belief that the
mathematics presented here will continue to influence the development of theoretical
economics, and have intensified my regret about not having completed it sooner.
There is a vast literature on this topic, which has influenced me in many ways,
and which cannot be described in any useful way here. Even so, I should say
something about how the present work stands in relation to three other books on
fixed points. Fixed Point Theorems with Applications to Economics and Game
Theory by Kim Border (1985) is a complement, not a substitute, explaining various
forms of the fixed point principle such as the KKMS theorem and some of the
many theorems of Ky Fan, along with the concrete details of how they are actually
applied in economic theory. Fixed Point Theory by Dugundji and Granas (2003) is,
even more than this book, a comprehensive treatment of the topic. Its fundamental
point of view (applications to nonlinear functional analysis) audience (professional
mathematicians) and technical base (there is extensive use of algebraic topology)
are quite different, but it is still a work with much to offer to economics. Particularly
notable is the extensive and meticulous information concerning the literature and
history of the subject, which is full of affection for the theory and its creators. The
book that was, by far, the most useful to me, is The Lefschetz Fixed Point Theorem
by Robert Brown (1971). Again, his approach and mine have differences rooted in
the nature of our audiences, and the overall objectives, but at their cores the two
books are quite similar, in large part because I borrowed a great deal.
I would like to thank the many people who, over the years, have commented
favorably on Selected Topics. It is a particular pleasure to acknowledge some very
detailed and generous written comments by Klaus Ritzberger. This work would not
have been possible without the support and affection of my families, both present
and past, for which I am forever grateful.
Contents
1 Introduction and Summary
1.1 The First Fixed Point Theorems
1.2 Fixing Kakutanis Theorem .
1.3 Essential Sets of Fixed Points .
1.4 Index and Degree . . . . . . . .
1.4.1 Manifolds . . . . . . . .
1.4.2 The Degree . . . . . . .
1.4.3 The Fixed Point Index .
1.5 Topological Consequences . . .
1.6 Dynamical Systems . . . . . . .
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Topological Methods
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2
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. . .
Sets
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23
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iii
CONTENTS
4.5.1
4.5.2
4.5.3
4.5.4
4.5.5
4.5.6
Continuity of Union . . . . . . . . .
Continuity of Intersection . . . . . .
Singletons . . . . . . . . . . . . . . .
Continuity of the Cartesian Product
The Action of a Function . . . . . . .
The Union of the Elements . . . . . .
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7 Retracts
7.1 Kinoshitas Example . . . . . . .
7.2 Retracts . . . . . . . . . . . . . .
7.3 Euclidean Neighborhood Retracts
7.4 Absolute Neighborhood Retracts
7.5 Absolute Retracts . . . . . . . . .
7.6 Domination . . . . . . . . . . . .
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II
Smooth Methods
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115
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10 Differentiable Manifolds
125
10.1 Review of Multivariate Calculus . . . . . . . . . . . . . . . . . . . . 126
10.2 Smooth Partitions of Unity . . . . . . . . . . . . . . . . . . . . . . . 128
CONTENTS
10.3
10.4
10.5
10.6
10.7
10.8
10.9
Manifolds . . . . . . . . . . . . . . . .
Smooth Maps . . . . . . . . . . . . . .
Tangent Vectors and Derivatives . . . .
Submanifolds . . . . . . . . . . . . . .
Tubular Neighborhoods . . . . . . . . .
Manifolds with Boundary . . . . . . .
Classification of Compact 1-Manifolds .
11 Sards Theorem
11.1 Sets of Measure Zero . . . . . . . .
11.2 A Weak Fubini Theorem . . . . . .
11.3 Sards Theorem . . . . . . . . . . .
11.4 Measure Zero Subsets of Manifolds
11.5 Genericity of Transversality . . . .
12 Degree Theory
12.1 Orientation . . . . . . . . . . . . .
12.2 Induced Orientation . . . . . . . .
12.3 The Degree . . . . . . . . . . . . .
12.4 Composition and Cartesian Product
13 The
13.1
13.2
13.3
13.4
13.5
III
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131
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14 Topological Consequences
14.1 Euler, Lefschetz, and Eilenberg-Montgomery
14.2 The Hopf Theorem . . . . . . . . . . . . . .
14.3 More on Maps Between Spheres . . . . . . .
14.4 Invariance of Domain . . . . . . . . . . . . .
14.5 Essential Sets Revisited . . . . . . . . . . . .
15 Vector Fields and their Equilibria
15.1 Euclidean Dynamical Systems . . .
15.2 Dynamics on a Manifold . . . . . .
15.3 The Vector Field Index . . . . . . .
15.4 Dynamic Stability . . . . . . . . . .
15.5 The Converse Lyapunov Problem .
15.6 A Necessary Condition for Stability
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Chapter 1
Introduction and Summary
The Brouwer fixed point theorem states that if C is a nonempty compact convex
subset of a Euclidean space and f : C C is continuous, then f has a fixed point,
which is to say that there is an x C such that f (x ) = x . The proof of
this by Brouwer (1912) was one of the major events in the history of topology.
Since then the study of such results, and the methods used to prove them, has
flourished, undergoing radical transformations, becoming increasingly general and
sophisticated, and extending its influence to diverse areas of mathematics.
Around 1950, most notably through the work of Nash (1950, 1951) on noncooperative games, and the work of Arrow and Debreu (1954) on general equilibrium
theory, it emerged that in economists most basic and general models, equilibria
are fixed points. The most obvious consequence of this is that fixed point theorems provide proofs that these models are not vacuous. But fixed point theory also
informs our understanding of many other issues such as comparative statics, robustness under perturbations, stability of equilibria with respect to dynamic adjustment
processes, and the algorithmics and complexity of equilibrium computation. In particular, since the mid 1970s the theory of games has been strongly influenced by
refinement concepts defined largely in terms of robustness with respect to certain
types of perturbations.
As the range and sophistication of economic modelling has increased, more advanced mathematical tools have become relevant. Unfortunately, the mathematical
literature on fixed points is largely inaccessible to economists, because it relies heavily on homology. This subject is part of the standard graduate school curriculum
for mathematicians, but for outsiders it is difficult to penetrate, due to its abstract
nature and the amount of material that must be absorbed at the beginning before
the structure, nature, and goals of the theory begin to come into view. Many researchers in economics learn advanced topics in mathematics as a side product of
their research, but unlike infinite dimensional analysis or continuous time stochastic
processes, algebraic topology will not gradually achieve popularity among economic
theorists through slow diffusion. Consequently economists have been, in effect,
shielded from some of the mathematics that is most relevant to their discipline.
This monograph presents an exposition of advanced material from the theory of
fixed points that is, in several ways, suitable for graduate students and researchers
in mathematical economics and related fields. In part the fit with the intended
2
1.1
theorem are closely related to algorithmic procedures for finding approximate fixed
points. Chapter 3 discusses the best known general algorithm due to Scarf, a new
algorithm due to the author and Rabee Tourky, and homotopy methods, which are
the most popular in practice, but require differentiability. The last decade has seen
major breakthroughs in computer science concerning the computational complexity
of computing fixed points, with particular reference to (seemingly) simple games
and general equilibrium models. These developments are sketched briefly in Section
3.7.
In economics and game theory fixed point theorems are most commonly used to
prove that a model has at least one equilibrium, where an equilibrium is a vector of
endogenous variable for the model with the property that each individual agents
predicted behavior is rational, or utility maximizing, if that agent regards all
the other endogenous variables as fixed. In economics it is natural, and in game
theory unavoidable, to consider models in which an agent might have more than
one rational choice. Our first generalization of Brouwers theorem addresses this
concern.
If X and Y are sets, a correspondence F : X Y is a function from X to the
nonempty subsets of Y . (On the rare occasions when they arise, we use the term set
valued mapping for a function from X to all the subsets of Y , including the empty
set.) We will tend to regard a function as a special type of correspondence, both
intuitively and in the technical sense that we will frequently blur the distinction
between a function f : X Y and the associated correspondence x 7 {f (x)}.
If Y is a topological space, F is compact valued if, for all x X, F (x) is
compact. Similarly, if Y is a subset of a vector space, then F is convex valued if
each F (x) is convex.
The extension of Brouwers theorem to correspondences requires a notion of
continuity for correspondences. If X and Y are topological spaces, a correspondence
F : X Y is upper semicontinuous if it is compact valued and, for each x0 X
and each neighborhood V Y of F (x0 ), there is a neighborhood U X of x0 such
that F (x) V for all x U. It turns out that if X and Y are metric spaces and Y
is compact, then F is upper semicontinuous if and only if its graph
Gr(F ) := { (x, y) X Y : y F (x) }
is closed. (Proving this is a suitable exercise, if you are so inclined.) Thinking of
upper semicontinuity as a matter of the graph being closed is quite natural, and in
economics this condition is commonly taken as definition, as in Debreu (1959). In
Chapter 5 we will develop a topology on the space of nonempty compact subsets
of Y such that F is upper semicontinuous if and only if it is a continuous function
relative to this topology.
A fixed point of a correspondence F : X X is a point x X such that
1.2
Mathematicians strive to craft theorems that maximize the strength of the conclusions while minimizing the strength of the assumptions. One reason for this is
obvious: a stronger theorem is a more useful theorem. More important, however, is
the desire to attain a proper understanding of the principle the theorem expresses,
and to achieve an expression of this principle that is unencumbered by useless clutter. When a theorem that is too weak is proved using methods that happen
to work there is a strong suspicion that attempts to improve the theorem will
uncover important new concepts. In the case of Brouwers theorem the conclusion,
that the space has the fixed point property, is a purely topological assertion. The
assumption that the space is convex, and in Kakutanis theorem the assumption
that the correspondences values are convex, are geometric conditions that seems
out of character and altogether too strong. Suitable generalizations were developed
after World War II.
A homotopy is a continuous function h : X [0, 1] Y where X and Y are
topological spaces. It is psychologically natural to think of the second variable in the
domain as representing time, and we let ht := h(, t) : X Y denote the function
at time t, so that h is a process that continuously deforms a function h0 into h1 .
Another intuitive picture is that h is a continuous path in the space C(X, Y ) of
continuous function from X to Y . As we will see in Chapter 5, this intuition can
be made completely precise: when X and Y are metric spaces and X is compact,
there is a topology on C(X, Y ) such that a continuous path h : [0, 1] C(X, Y ) is
the same thing as a homotopy.
We say that two functions f, g : X Y are homotopic if there is a homotopy
h with h0 = f and h1 = g. This is easily seen to be an equivalence relation:
symmetry and reflexivity are obvious, and to establish transitivity we observe that
if e is homotopic to f and f is homotopic to g, then there is a homotopy between
e and g that follows a homotopy between e and f at twice its original speed, then
follows a homotopy between f and g at double the pace. The equivalence classes
are called homotopy classes.
A space X is contractible if the identity function IdX is homotopic to a constant
function. That is, there is a homotopy c : X [0, 1] X such that c0 = IdX and
c1 (X) is a singleton; such a homotopy is called a contraction. Convex sets are
contractible. More generally, a subset X of a vector space is star-shaped if there
is x X (the star) such that X contains the line segment
{ (1 t)x + tx : 0 t 1 }
between each x X and x . If X is star-shaped, there is a contraction
(x, t) 7 (1 t)x + tx .
It seems natural to guess that a nonempty compact contractible space has the
fixed point property. Whether this is the case was an open problem for several years,
but it turns out to be false. In Chapter 7 we will see an example due to Kinoshita
(1953) of a nonempty compact contractible subset of R3 that does not have the
fixed point property. Fixed point theory requires some additional ingredient.
If X is a topological space, a subset A X is a retract if there is a continuous
function r : X A with r(a) = a for all a A. Here we tend to think of X as a
simple space, and the hope is that although A might seem to be more complex,
or perhaps crumpled up, it nonetheless inherits enough of the simplicity of X. A
particularly important manifestation of this is that if r : X A is a retraction and
X has the fixed point property, then so does A, because if f : A A is continuous,
then so is f r : X A X, so f r has a fixed point, and this fixed point
necessarily lies in A and is consequently a fixed point of f . Also, a retract of a
contractible space is contractible because if c : X [0, 1] X is a contraction of
X and r : X A X is a retraction, then
(a, t) 7 r(c(a, t))
is a contraction of A.
A set A Rm is a Euclidean neighborhood retract (ENR) if there is an
open superset U Rm of A and a retraction r : U A. If X and Y are metric
spaces, an embedding of X in Y is a function e : X Y that is a homeomorphism
between X and e(X). That is, e is a continuous injection1 whose inverse is also
continuous when e(X) has the subspace topology inherited from Y . An absolute
neighborhood retract (ANR) is a separable2 metric space X such that whenever
Y is a separable metric space and e : X Y is an embedding, there is an open
superset U Y of e(X) and a retraction r : U e(X). This definition probably
seems completely unexpected, and its difficult to get any feeling for it right away.
In Chapter 7 well see that ANRs have a simple characterization, and that many
of the types of spaces that come up most naturally are ANRs, so this condition
is quite a bit less demanding than one might guess at first sight. In particular, it
will turn out that every ENR is an ANR, so that being an ENR is an intrinsic
property insofar as it depends on the topology of the space and not on how the
space is embedded in a Euclidean space.
An absolute retract (AR) is a separable metric space X such that whenever
Y is a separable metric space and e : X Y is an embedding, there is a retraction
r : Y e(X). In Chapter 7 we will prove that an ANR is an AR if and only if it
is contractible.
Theorem 1.2.1. If C is a nonempty compact AR and F : C C is an upper
semicontinuous contractible valued correspondence, then F has a fixed point.
An important point is that the values of F are not required to be ANRs.
1
We will usually use the terms injective rather than one-to-one, surjective rather than
onto, and bijective to indicate that a function is both injective and surjective. An injection is
an injective function, a surjection is a surjective function, and a bijection is a bijective function.
2
A metric space is separable if it has a countable dense subset.
For practical purposes this is the maximally general topological fixed point
theorem, but for mathematicians there is an additional refinement. There is a concept called acyclicity that is defined in terms of the concepts of algebraic topology.
A contractible set is necessarily acyclic, but there are acyclic spaces (including compact ones) that are not contractible. The famous Eilenberg-Montgomery fixed point
theorem is:
Theorem 1.2.2 (Eilenberg and Montgomery (1946)). If C is a nonempty compact
AR and F : C C is an upper semicontinuous acyclic valued correspondence, then
F has a fixed point.
1.3
It might seem like we have already reached a satisfactory and fitting resolution
of The Fixed Point Problem, but actually (both in pure mathematics and in
economics) this is just the beginning. You see, fixed points come in different flavors.
1
b
Figure 1.1
The figure above shows a function f : [0, 1] [0, 1] with two fixed points, s
and t. If we perturb the function slightly by adding a small positive constant, s
disappears in the sense that the perturbed function does not have a fixed point
anywhere near s, but a function close to f has a fixed point near t. More precisely,
if X is a topological space and f : X X is continuous, a fixed point x of f is
essential if, for any neighborhood U of x , there is a neighborhood V of the graph
of f such that any continuous f : X X whose graph is contained in V has a
fixed point in U. If a fixed point is not essential, then we say that it is inessential.
These concepts were introduced by Fort (1950).
There need not be an essential fixed point. The function shown in Figure 1.2
has an interval of fixed points. If we shift the function down, there will be a fixed
point near the lower endpoint of this interval, and if we shift the function up there
will be a fixed point near the upper endpoint.
This example suggests that we might do better to work with sets of fixed points.
A set S of fixed points of a function f : X X is essential if it is closed, it has a
neighborhood that contains no other fixed points, and for any neighborhood U of S,
there is a neighborhood V of the graph of f such that any continuous f : X X
whose graph is contained in V has a fixed point in U. The problem with this
concept is that large connected sets are not of much use. For example, if X is
compact and has the fixed point property, then the set of all fixed points of f is
essential. It seems that we should really be interested in sets of fixed points that are
either essential and connected3 or essential and minimal in the sense of not having
a proper subset that is also essential.
1
1
Figure 1.2
In Chapter 8 we will show that any essential set of fixed points contains a minimal essential set, and that minimal essential sets are connected. The theory of
refinements of Nash equilibrium (e.g., Selten (1975); Myerson (1978); Kreps and
Wilson (1982); Kohlberg and Mertens (1986); Mertens (1989, 1991); Govindan and
Wilson (2008)) has many concepts that amount to a weakening of the notion of
essential set, insofar as the set is required to be robust with respect to only certain types of perturbations of the function or correspondence. In particular, Jiang
(1963) pioneered the application of the concept to game theory, defining an essential!Nash equilibrium and an essential set of Nash equilibria in terms
of robustness with respect to perturbations of the best response correspondence
induced by perturbations of the payoffs. The mathematical foundations of such
3
We recall that a subset S of a topological space X is connected if there do not exist two
disjoint open sets U1 and U2 with S U1 6= 6= S U2 and S U1 U2 .
1.4
There are different types of essential fixed points. Figure 1.3 shows a function
with three fixed points. At two of them the function starts above the diagonal and
goes below it as one goes from left to right, and at the third it is the other way
around. For any k it is easy to imagine a function with k fixed points of the first
type and k 1 fixed points of the second type.
This phenomenon generalizes to higher dimensions. Let
D m = { x Rm : kxk 1 } and S m1 = { x Rm : kxk = 1 }
be the m-dimensional unit disk and the (m 1)-dimensional unit sphere, and suppose that f : D m D m is a C function. In the best behaved case each fixed
point x is in the interior D m \ S m1 of the disk and regular, which means that
IdRm Df (x ) is nonsingular, where Df (x ) : Rm Rm is the derivative of f at x .
We define the index of x to be 1 if the determinant of IdRm Df (x ) is positive
and 1 if this determinant is negative. We will see that there is always one more
fixed point of index 1 than there are fixed points of index 1, which is to say that
the sum of the indices is 1.
What about fixed points on the boundary of the disk, or fixed points that arent
regular, or nontrivial connected sets of fixed points? What about correspondences?
What happens if the domain is a possibly infinite dimensional ANR? The most challenging and significant aspect of our work will be the development of an axiomatic
theory of the index that is general enough to encompass all these possibilities. The
work proceeds through several stages, and we describe them in some detail now.
1
b
1
Figure 1.3
10
1.4.1
Manifolds
11
The inverse and implicit function theorems have important generalizations. The
point p is a regular point of f if the image of Df (p) is all of Tf (p) N. We say that
f : M N is a C r diffeomorphism if m = n, f is a bijection, and both f and f 1
are C r . The generalized inverse function theorem asserts that if m = n, f : M N
is C r , and p is a regular point of f , then there is an open U M containing p such
that f (U) is an open subset of N and f |U : U f (U) is a C r diffeomorphism.
If 0 s m, a set S Rk is an s-dimensional C r submanifold of M if it
is an s-dimensional C r submanifold that happens to be contained in M. We say
that q N is a regular value of f if every p f 1 (q) is a regular point. The
generalized implicit function theorem, which is known as the regular value theorem,
asserts that if q is a regular value of f , then f 1 (q) is an (m n)-dimensional C r
submanifold of M.
1.4.2
The Degree
The degree is closely related to the fixed point index, but it has its own theory,
which has independent interest and significance. The approach we take here is to
work with the degree up to the point where its theory is more or less complete, then
translate what we have learned into the language of the fixed point index.
We now need to introduce the concept of orientation. Two ordered bases
v1 , . . . , vm and w1 , . . . , wm of an m-dimensional vector space have the same orientation if the determinant of the linear transformation taking each vi to wi is
positive. It is easy to see that this is an equivalence relation with two equivalence
classes. An oriented vector space is a finite dimensional vector space with a designated orientation whose elements are said to be positively oriented. If V and
W are m-dimensional oriented vector spaces, a nonsingular linear transformation
L : V W is orientation preserving if it maps positively oriented ordered bases
of V to positively oriented ordered bases of W , and otherwise it is orientation
reversing. For an intuitive appreciation of this concept just look in a mirror: the
linear map taking each point in the actual world to its position as seen in the mirror
is orientation reversing, with right shoes turning into left shoes and such.
In our discussion of degree theory nothing is lost by working with C objects
rather than C r objects for general r, and smooth will be a synonym for C . An
orientation for a smooth manifold M is a continuous specification of an orientation
of each of the tangent spaces Tp M. We say that M is orientable if it has an
orientation; the most famous examples of unorientable manifolds are the Mobius
strip and the Klein bottle. (From a mathematical point of view 2-dimensional
projective space is perhaps more fundamental, but it is difficult to visualize.) An
oriented manifold is a manifold together with a designated orientation.
If M and N are oriented smooth manifolds of the same dimension, f : M N
is a smooth map, and p is a regular point of f , we say that f is orientation
preserving at p if Df (p) : Tp M Tf (p) N is orientation preserving, and otherwise
f is orientation reversing at p. If q is a regular value of f and f 1 (q) is finite,
then the degree of f over q, denoted by deg
q (f ), is the number of points in
1
f (q) at which f is orientation preserving minus the number of points in f 1 (q)
at which f is orientation reversing.
12
We need to extend the degree to situations in which the target point q is not a
regular value of f , and to functions that are merely continuous. Instead of being
able to define the degree directly, as we did above, we will need to proceed indirectly,
showing that the generalized degree is determined by certain of its properties, which
we treat as axioms.
The first step is to extend the concept, giving it a local character. For a
compact C M let C = C (M \ C) be the topological boundary of C, and let
int C = C \ C be its interior. A smooth function f : C N with compact domain
C M is said to be smoothly degree admissible over q N if f 1 (q) C =
and q is a regular value of f . As above, for such a pair (f, q) we define deg
q (f ) to
be the number of p f 1 (q) at which f is orientation preserving minus the number
1
whenever C is a compact subset of C and f (q) has an empty intersection with
the closure of C \ C . Also, if C = C1 C2 where C1 and C2 are compact and
disjoint, then
deg
q (f ) = deg q (f |C1 ) + deg q (f |C2 ).
From the point of view of topology, what makes the degree important is its
invariance under homotopy. If C M is compact, a smooth homotopy h : C
[0, 1] N is smoothly degree admissible over q if h1 (q) (C [0, 1]) =
and q is a regular value of h0 and h1 . In this circumstance
deg
q (h0 ) = deg q (h1 ).
()
1
b
+1
b
+1
b
+1
1
b
+1
b
t=0
t=1
Figure 1.4
13
deg
q (h0 ) = deg q (h0 ) and deg q (h1 ) = deg q (h1 ) when q is sufficiently close to q.
It turns out that the smooth degree is completely characterized by the properties
we have seen. That is, if D (M, N) is the set of pairs (f, q) in which f : C N
is smoothly degree admissible over q, then (f, q) 7 deg
q (f ) is the unique function
1
(q) is a singleton {p}
(1) deg
q (f ) = 1 for all (f, q) D (M, N) such that f
14
(2) deg
q (f ) =
i=1 deg q (f |Ci ) whenever (f, q) D (M, N), the domain of f is
C, and C1 , . . . , Cr are pairwise disjoint compact subsets of C such that
f 1 (q) int C1 . . . int Cr .
(3) deg
q (h0 ) = deg q (h1 ) whenever C M is compact and the homotopy h :
C [0, 1] N is smoothly degree admissible over q.
We note two additional properties of the smooth degree. The first is that if, in
addition to M and N, M and N are m -dimensional smooth functions, (f, q)
D (M, N), and (f , q ) D (M , N ), then
(f f , (q, q )) D (M M , N N ) and
deg
(q,q ) (f f ) = degq (f )degq (f ).
deg
q (f ) = deg q (f )
15
1.4.3
Although the degree can be applied to continuous functions, and even to convex
valued correspondences, it is restricted to finite dimensional manifolds. For such
spaces the fixed point index is merely a reformulation of the degree. Its application
to general equilibrium theory was initiated by Dierker (1972), and it figures in the
analysis of the Lemke-Howson algorithm of Shapley (1974). There is also a third
variant of the underlying principle, for vector fields, that is developed in Chapter 15,
and which is related to the theory of dynamical systems. Hofbauer (1990) applied
the vector field index to dynamic issues in evolutionary stability, and Ritzberger
(1994) applies it systematically to normal form game theory.
However, it turns out that the fixed point index can be generalized much further,
due to the fact that, when we are discussing fixed points, the domain and the
range are the same. The general index is developed in three main stages. In
order to encompass these stages in a single system of terminology and notation we
take a rather abstract approach. Fix a metric space X. An index admissible
correspondence for X is an upper semicontinuous correspondence F : C X,
where C X is compact, that has no fixed points in C. An index base for X is
a set I of index admissible correspondences such that:
(a) f I whenever C X is compact and f : C X is an index admissible
continuous function;
(b) F |D I whenever F : C X is an element of I, D C is compact, and
F |D is index admissible.
Definition 1.4.1. Let I be an index base for X. An index for I is a function
X : I Z satisfying:
(I1) (Normalization) If c : C X is a constant function whose value is an element
of int C, then X (c) = 1.
(I2) (Additivity) If F : C X is an element of I, C1 , . . . , Cr are pairwise disjoint
compact subsets of C, and F P(F ) int C1 . . . int Cr , then
X
X (F ) =
X (F |Ci ).
i
16
g g| ) IS (X),
with X, X SS , (D, g g|D ) IS (X), and (D,
D
X (
g g|D ) = X (g g|D ).
The index is said to be multiplicative if:
(M) (Multiplication) If X, X SS , F IS (X), and F IS (X ), then
XX (F F ) = X (F ) X (F ).
17
Let SS Ctr be the class of ANRs, and for each X SS Ctr let IS Ctr (X) be the union over
compact C X of the sets of index admissible upper semicontinuous contractible
valued correspondences F : C X. The central goal of this book is:
Theorem 1.4.4. There is a unique index Ctr for S Ctr , which is multiplicative.
The passage from the indices Rm to Ctr has two stages. The first exploits
Commutativity to extend from Euclidean spaces and continuous functions to ANRs
and continuous functions. There is a significant result that is the technical basis for
this. Let X be a metric space with metric d. If Y is a topological space and > 0,
a homotopy : Y [0, 1] X is an -homotopy if
d (y, s), (y, t) <
for all y Y and all 0 s, t 1. We say that h0 and h1 are -homotopic. For
> 0, a topological space D -dominates C X if there are continuous functions
: C D and : D X such that : C X is -homotopic to IdC . In
Section 7.6 we show that:
1.5
Topological Consequences
The final section of the book develops applications of the index. Chapter 14
presents a number of classical concepts and results from topology that are usually
proved homologically. Let X be a compact ANR. The Euler characteristic of X
is the index of IdX . If F : X X is an upper semicontinuous contractible valued
correspondence, the index of F is called the Lefschetz number of F . Of course
Additivity implies that F has a fixed point if its Lefschetz number is not zero.
The celebrated Lefschetz fixed point theorem is this assertion (usually restricted to
18
compact manifolds and continuous functions) together with a homological characterization of the Lefschetz number. If X is contractible, then the Lefschetz number
of any F : X X is equal to the Euler characteristic of F , which is one. Thus
we arrive at our version of the Eilenberg-Montgomery theorem: if X is a compact
AR and F : X X is a upper semicontinuous contractible valued correspondence,
then F has a fixed point.
Chapter 14 also develops many of the classical theorems concerning maps between spheres. The most basic of these is Hopfs theorem: two continuous functions
f, f : S m S m are homotopic if and only if they have the same degree, so that the
degree is a complete homotopy invariant for maps between spheres of the same
dimension. There are many other theorems concerning maps between spheres of
the same dimension. Of these, one in particular has greater depth: if f : S m S m
is continuous and f (p) = f (p) for all p S m , then the degree of f is odd. This
and its many corollaries constitute the Borsuk-Ulam theorem.
Using these results, we prove the frequently useful theorem known as invariance
of domain: if U Rm is open and f : U Rm is continuous and injective, then
f (U) is open and f is a homeomorphism onto its image.
If a connected set of fixed points has nonzero index, then it is essential, by virtue
of Continuity. The result in Section 14.5 shows that the converse holds for convex
valued correspondences with convex domains, so for the settings most commonly
considered in economics the notion of essentiality does not have independent significance. But it is important to understand that this result does not imply that
a component of the set of Nash equilibria of a normal form game of index zero is
inessential in the sense of Jiang (1963). In fact Hauk and Hurkens (2002) provide
a concrete example of an essential component of index zero.
1.6
Dynamical Systems
Dynamic stability is a problematic issue for economic theory. On the one hand,
particularly in complex settings, it seems that an equilibrium cannot a plausible
prediction unless it can be understood as the end state of a dynamic adjustment
process for which it is dynamically stable. In physics and chemistry there are explicit dynamical systems, and with respect to those stability is a well accepted principle. But in economics, explicit models of dynamic adjustment are systematically
inconsistent with the principle of rational expectations: if a model of continuous
adjustment of prices, or of mixed strategies, is understood and anticipated by the
agents in the model, their behavior will exploit the process, not conform to it.
Early work in general equilibrium theory (e.g., Arrow and Hurwicz (1958); Arrow et al. (1959)) found special cases, such as a single agent or two goods, in which
at least one equilibrium is necessarily stable with respect to natural price adjustment processes. But Scarf (1960) produced examples showing that one could not
hope for more general positive results, in the sense that naive dynamic adjustment
processes, such as Walrasian tatonnement, can easily fail to have stable dynamics,
even when there is a unique equilibrium and as few as three goods. A later stream
of research (Saari and Simon (1978); Saari (1985); Williams (1985); Jordan (1987))
19
showed that stability is informationally demanding, in the sense that an adjustment process that is guaranteed to return to equilibrium after a small perturbation
requires essentially all the information in the matrix of partial derivatives of the
aggregate excess demand function. On the whole there seems to be little hope of
finding a theoretical basis for an assertion that some equilibrium is stable, or that
a stable equilibrium exists.
In his Foundations of Economic Analysis Samuelson (1947) Samuelson describes
a correspondence principle, according to which the stability of an equilibrium
has implications for the qualitative properties of its comparative statics. In this
style of reasoning the stability of a given equilibrium is a hypothesis rather than a
conclusion, so the problematic state of the existence issue is less relevant. That is,
instead of claiming that some dynamical process should result in a stable equilibrium, one argues that equilibria with certain properties are not stable, so if what
we observe is an equilibrium, it cannot have these properties.
Proponents of such reasoning still need to wrestle with the fact that there is no
canonical dynamical process. (The conceptual foundations of economic dynamics,
and in particular the principle of rational expectations, were not well understood in
Samuelsons time, and his discussion would be judged today to have various weaknesses.) Here there is the possibility of arguing that although any one dynamical
process might be ad hoc, the instability is common to all reasonable or natural
dynamics, for example those in which price adjustment is positively related to excess
demand, or that each agents mixed strategy adjusts in a direction that would improve her expected utility if other mixed strategies were not also adjusting. From a
strictly logical point of view, such reasoning might seem suspect, but it seems quite
likely that most economists find it intuitively and practically compelling.
In Chapter 15 we present a necessary condition for stability of a component
of the set of equilibria that was introduced into game theory by Demichelis and
Ritzberger (2003). (See also Demichelis and Germano (2000).) We now give an
informal description of this result, with the relevant background, and relate it to
Samuelsons correspondence principle.
Let M be an m-dimensional C 2 manifold, where r 2. A vector field on a
set S M is a continuous (in the obvious sense) assignment of a tangent vector
p Tp M to each p S. Vector fields have many applications, but by far the most
important is that if is defined on an open U M and satisfies a mild technical
condition, then it determines an autonomous dynamical system: there is an open
W U R such that for each p U, { t R : (p, t) W } is an interval containing
0, and a unique function : W U such that (p, 0) = p for all p and, for each
(p, t) W , the time derivative of at (p, t) is (p,t) . If W is the maximal domain
admitting such a function, then is the flow of . A point p where p = 0 is an
equilibrium of .
A set A M is invariant if (p, t) A for all p A and t 0. The -limit
set of p M is
\
{ (p, t) : t t0 }.
t0 0
20
d
f ((p, t))|t=0
dt
21
for small positive t the flow (, t) will map neighborhoods of A into themselves.
The Lyapunov function given by the converse Lyapunov theorem is used in Section
15.6 to show that if A is dynamically stable and an ANR (otherwise the Euler characteristic is undefined) then the vector field index of A is (1)m (A). In particular,
if A is a singleton, then A can only be stable when the vector field index of A is
(1)m . This is the result of Demichelis and Ritzberger. The special case when
A = {p0 } is a singleton is a prominent result in the theory of dynamical systems
due to Krasnoselski and Zabreiko (1984).
We now describe the relationship between this result and qualitative properties
of an equilibriums comparative statics. Consider the following stylized example.
Let U be an open subset of Rm ; an element of U is thought of as a vector of
endogenous variables. Let P be an open subset of Rn ; an element of P is thought of
as a vector of exogenous parameters. Let z : U P Rm be a C 1 function, and let
x z(x, ) and z(x, ) denote the matrices of partial derivatives of the components
of z with respect to the components of x and .
We think of z as a parameterized vector field on U. An equilibrium for a
parameter P is an x U such that z(x, ) = 0. Suppose that x0 is an
equilibrium for 0 , and x z(x0 , 0 ) is nonsingular. The implicit function theorem
gives a neighborhood V of and a C 1 function : V U with (0 ) = x0 and
z((), ) = 0 for all V . The method of comparative statics is to differentiate
this equation with respect to at 0 , then rearrange, obtaining the equation
d
(0 ) = x z(x0 , 0 )1 (x0 , 0 )
d
describing how the endogenous variables adjust, in equilibrium, to changes in the
vector of parameters. The Krasnoselski-Zabreiko theorem implies that if {x0 } is an
asymptotically stable set for the dynamical system determined by the vector field
z(, 0 ), then the determinant of x z(x0 , 0 )1 is positive. This is a precise and
general statement of the correspondence principle.
Part I
Topological Methods
22
Chapter 2
Planes, Polyhedra, and Polytopes
This chapter studies basic geometric objects defined by linear equations and
inequalities. This serves two purposes, the first of which is simply to introduce
basic vocabulary. Beginning with affine subspaces and half spaces, we will proceed to (closed) cones, polyhedra, and polytopes, which are polyhedra that are
bounded. A rich class of well behaved spaces is obtained by combining polyhedra
to form polyhedral complexes. Although this is foundational, there are nonetheless
several interesting and very useful results and techniques, notably the separating
hyperplane theorem, Farkas lemma, and barycentric subdivision.
2.1
Affine Subspaces
Throughout the rest of this chapter we work with a fixed d-dimensional real inner
product space V . (Of course we are really talking about Rd , but a more abstract
setting emphasizes the geometric nature of the constructions and arguments.) We
assume familiarity with the concepts and results of basic linear algebra.
An affine combination of y0 , . . . , yr V is a point of the form
0 y 0 + + r y r
where = (0 , . . . , r ) is a vector of real numbers whose components sum to 1. We
say that y0 , . . . , yr are affinely dependent if it is possible to represent a point as
an affine combination of these points in two different ways: that is, if
X
X
X
X
j = 1 =
j and
j y j =
j yj ,
j
then = . If y0 , . . . , yr are not affinely dependent, then they are affinely independent.
Lemma 2.1.1. For any y0 , . . . , yr V the following are equivalent:
(a) y0 , . . . , yr are affinely independent;
(b) y1 y0 , . . . , yr y0 are linearly independent;
23
24
(c) therePdo not exist 0 , . . . , r R, not all of which are zero, with
and j j yj = 0.
j = 0
If P
we set j = j
Pj , then j j = 0 and j j yj = 0, so (c) implies (a). In turn,
if j j = 0 and j j yj = 0, then
1 (y1 y0 ) + + r (yr y0 ) = (1 + + r )y0 + 1 y1 + + r yr = 0,
2.2
25
26
yj x
hx, ni
,n
0.
kyj xk
kyj xk
The lineality space of a convex set C is
LC = RC RC = { y V : x + y C for all x C and R }.
The lineality space is closed under addition and scalar multiplication, so it is a
linear subspace of V , and in fact it is the largest linear subspace of V contained in
RC . Let L
C be the orthogonal complement of LC . Clearly C + LC = C, so
C = (C L
C ) + LC .
A convex cone is said to be pointed if its lineality space is {0}.
Lemma 2.2.4. If C =
6 V is a closed convex cone, then there is n C with hn, xi > 0
for all x C \ LC .
Proof. For n C let Zn = { x C : hx, ni = 0 }. Let n be a point in C that
minimizes the dimension of the span of Zn . Aiming at a contradiction, suppose
that 0 6= x Zn \ LC . Then x
/ C because x
/ LC , and Farkas Lemma gives
an n C with hx, n i < 0. Then Zn+n Zn Zn (this inclusion holds for all
n, n C ) and the span of Zn+n does not contain x, so it is a proper subspace of
the span of Zn .
2.3
Polyhedra
27
2.3. POLYHEDRA
28
X
X
x,
j nj
j j
jJ
jJ
I
for
some
J,
that
F
=
F
i
iJ
iJ Ii
T
for some J , and that F = P iJJ Ii is a face of P .
Corollary 2.3.7. The facets of P are F{1} , . . . , F{k} . The dimension of each F{i}
is one less than the dimension of P , The facets are the only faces of P with this
dimension.
Proof. Minimality implies that each F{i} is a proper face, and the result above
implies that F{i} cannot be a proper subset of another proper face. Thus each F{i}
is a facet.
For each i minimality implies that for each j 6= i there is some xj F{i} \ F{j} .
Let x be a convex combination of these with positive weights, then F{i} contains a
neighborhood of x in Ii , so the dimension of F{i} is the dimension of G Ii , which
is one less than the dimension of P .
A face F that is not a facet is a proper face of some facet, so its dimension is
not greater than two less than the dimension of P .
29
2.4. POLYTOPES
Now suppose that P is bounded. Any point in P that is not a vertex can be
written as a convex combination of points in proper faces of P . Induction on the
dimension of P proves that:
Proposition 2.3.8. If P is bounded, then it is the convex hull of its set of vertices.
An extreme point of a convex set is a point that is not a convex combination
of other points in the set. This result immediately implies that only vertices of P
can be extreme. In fact any vertex v is extreme: if {v} = P I where I is the
bounding hyperplane of a half space H containing P , then v cannot be a convex
combination of elements of P \ I.
2.4
Polytopes
30
Proof. Let L be its lineality, and let K be a linear subspace of V that is complementary to L in the sense that K L = {0} and K + L = V . Let Q = P K. Then
P = Q + L, and the lineality of Q is {0}, so RQ is pointed. Let S be the convex
hull of the set of initial points of Q. Above we saw that this is the convex hull of
the set of vertices of Q, so S is a polytope. Now Proposition 2.3.2 gives
P = L + RQ + S.
2.5
Polyhedral Complexes
A wide variety of spaces can be created by taking the union of a finite collection
of polyhedra.
Definition 2.5.1. A polyhedral complex is a finite set P = {P1 , . . . , Pk } of
polyhedra in V such that:
(a) F P whenever P P and F is a nonempty face of P ;
(b) for any 1 i, j k, Pi Pj is a common (possibly empty) face of Pi and Pj .
The underlying space of the complex is
|P| :=
P,
P P
31
b
b
b
b
b
b
b
b
b
b
b
b
b
b
32
This construction shows that the underlying space of a polytopal complex is also
the underlying space of a simplicial complex. In addition, repeating this process
can give a triangulation with small simplices. The diameter of a polytope is the
maximum distance between any two of its points. The mesh of a polytopal complex
is the maximum of the diameters of its polytopes.
Consider an -dimensional simplex P whose vertices are v0 , . . . , v . The barycenter of P is
1
(v0 + + v ).
(P ) :=
+1
In the construction above, suppose that P is a simplicial complex, and that we
chose wP = P for all P . We would like to bound the diameter of the simplices in
the subdivision of |P|, which amounts to giving a bound on the maximum distance
between the barycenters of any two nested faces. After reindexing, these can be
taken to be the faces spanned by v0 , . . . , vk and v0 , . . . , v where 0 k < m and
m is the dimension of P. The following rather crude inequality is sufficient for our
purposes.
1
1
(v0 + + vk )
(v0 + + v )
k+1
+1
X X
1
vi vj
=
(k + 1)( + 1) 0ik 0j
X
X
1
kvi vj k
(k + 1)( + 1)
0ik 0j,j6=i
1
m
(k + 1)D
D.
(k + 1)( + 1)
m+1
It follows from this that the mesh of the subdivision of |P| is not greater than
m/(m + 1) times the mesh of P. Since we can subdivide repeatedly:
Proposition 2.5.2. The underlying space of a polytopal complex has triangulations
of arbitrarily small mesh.
Simplicial complexes can be understood in purely combinatoric terms. An abstract simplicial complex is a pair (V, ) where V is a finite set of vertices and
is a collection of subsets of V with the property that whenever and
. The geometric interpretation is as follows. Let { ev : v V } be the standard unit basis vectors of RV : the v-component of ev is 1 and all other coordinates
are 0. (Probably most authors would work with R|V | , but our approach is simpler
and formally correct insofar as Y X is the set of functions from X to Y .) For each
nonempty let P be the convex hull of { ev : v }, and let P = . The
simplicial complex
P(V,) = { P : }
is called the canonical realization of (V, ).
Let P be a simplicial complex, and let V be the set of vertices of P. For each
P P let P = P V be the set of vertices of P , and let = { P : P P }.
It is easy to see that extending the map v 7 e affinely on each simplex induces
2.6. GRAPHS
33
2.6
Graphs
34
maximal if it not contained (in the obvious sense) in a longer path. Two vertices
are connected if they are the endpoints of a path. This is an equivalence relation,
and a component of G is one of the graphs consisting of an equivalence class and
the edges in G joining its vertices. We say that G is connected if it has only one
component, so that any two vertices are connected. A walk v0 v1 vr is a cycle if
r 3, v0 , . . . , vr1 are distinct, and vr = v0 . If G has no cycles, then it is said to
be acyclic. A connected acyclic graph is a tree.
The following simple fact is the only result from graph theory applied in this
book. It is sufficiently obvious that there would be little point in including a proof.
Lemma 2.6.1. If the degree of each of the vertices of G is at most two, then the
components of G are maximal paths, cycles, and vertices with no neighbors.
This simple principle underlies all the algorithms described in Chapter 3. There
are an even number of endpoints of paths in G. If it is known that an odd number
represent or embody a situation that is not what we are looking for, then the rest
do embody what we are looking for, and in particular the number of solutions is
odd, hence positive. If it is known that exactly one endpoint embodies what we are
not looking for, and that endpoint is easily computed, then we can find a solution
by beginning at that point and following the path to its other endpoint.
Chapter 3
Computing Fixed Points
When it was originally proved, Brouwers fixed point theorem was a major breakthrough, providing a resolution of several outstanding problems in topology. Since
that time the development of mathematical infrastructure has provided access to
various useful techniques, and a number of easier demonstrations have emerged, but
there are no proofs that are truly simple.
There is an important reason for this. The most common method of proving
that some mathematical object exists is to provide an algorithm that constructs it,
or some proxy such as an arbitrarily accurate approximation, but for fixed points
this is problematic. Naively, one might imagine a computational strategy that
tried to find an approximate fixed point by examining the value of the function at
various points, eventually halting with a declaration that a certain point was a good
approximation of a fixed point. For a function f : [0, 1] [0, 1] such a strategy
is feasible because if f (x) > x and f (x ) < x (as is the case if x = 0 and x = 1
unless one of these is a fixed point) then the intermediate value function implies
that there is a fixed point between x and x . According to the sign of f (x ) x ,
where x = (x+ x )/2, we can replace x or x with x , obtaining an interval with the
same property and half the length. Iterating this procedure provides an arbitrarily
fine approximation of a fixed point.
In higher dimensions such a computational strategy can never provide a guarantee that the output is actually near a fixed point. To say precisely what we mean
by this we need to be a bit more precise. Suppose you set out in search of a fixed
point of a continuous function f : X X (where X is nonempty, compact, and
convex subset of a Euclidean space) armed with nothing more than an oracle that
evaluates f . That is, the only computational resources you can access are the theoretical knowledge that f is continuous, and a black box that tells you the value of
f at any point in its domain that you submit to it. An algorithm is, by definition,
a computational procedure that is guaranteed to halt eventually, so our supposed
algorithm for computing a fixed point necessarily halts after sampling the oracle
finitely many times, say at x1 , . . . , xn , with some declaration that such-and-such is
at least an approximation of a fixed point. Provided that the dimension of X is
at least two, the Devil could now change the function to one that agrees with the
original function at every point that was sampled, is continuous, and has no fixed
points anywhere near the point designated by the algorithm. (One way to do this is
35
36
to replace f with h1 f h where h : X X is a suitable homeomorphism satisfying h(xi ) = xi and h(f (xi )) = f (xi ) for all i = 1, . . . , n.) The algorithm necessarily
processes the new function in the same way, arriving at the same conclusion, but
for the new function that conclusion is erroneous.
Our strategy for proving Brouwers fixed point theorem will, of necessity, be a
bit indirect. We will prove the existence of objects that we will describe as points
that are approximately fixed. (The exact nature of such objects will vary from
one proof to the next.) An infinite sequence of such points, with the error of
the approximation converging to zero, will have the property that each of its limit
points is a fixed point.
The proof that any sequence in a compact space has an accumulation point uses
the axiom of choice, and in fact Brouwers fixed point theorem cannot be proved
without it. The axiom of choice was rather controversial when it emerged, with
constructivists (Brouwer himself became one late in life) arguing that mathematics
should only consider objects whose definitions are, in effect, algorithms for computing the object in question, or at least a succession of finer and finer approximations.
It turns out that this is quite restrictive, so the should of the last sentence becomes quite puritanical, at least in comparison with the rich mathematics allowed
by a broader set of allowed definitions and accepted axioms, and constructivism has
almost completely faded out in recent decades.
This chapter studies two algorithmic ideas for computing points that are approximate fixed. One of these uses an algorithm for computing a Nash equilibrium
of a two person game. The second may be viewed as a matter of approximating the
given function or correspondence with an approximation that is piecewise linear in
the sense that its graph is a polyhedral complex. In both cases the algorithm traverses a path of edges in a polyhedral complex, and in the final section we explain
recent advances in computer science concerning such algorithms and the problems
they solve.
3.1
In a two person game each of the two players is required to choose an element
from a set of strategies, without being informed of the other players choice, and each
players payoff depends jointly on the pair of strategies chosen. A pair consisting
of a strategy for each agent is a Nash equilibrium if neither agent can do better by
switching to some other strategy. The mixed extension is the derived two person
game with the same two players in which each players set of strategies is the set of
probability measures on that players set of strategies in the original game. Payoffs
in the mixed extension are computed by taking expectations.
In a sense, our primary concern in this section and the next is to show that when
the sets of strategies in the given game are finite, the mixed extension necessarily has
a Nash equilibrium. But we will actually do something quite a bit more interesting
and significant, by providing an algorithm that computes a Nash equilibrium. We
will soon see that the existence result is a special case of the Kakutani fixed point
theorem. But actually this case is not so special because we will eventually
37
see that two person games can be used to approximate quite general fixed point
problems.
Formally, a finite two person game consists of:
(a) nonempty finite sets S = {s1 , . . . , sm } and T = {t1 , . . . , tn } of pure strategies for the two agents, who will be called agent 1 and agent 2;
(b) payoff functions u, v : S T R.
Elements of S T are called pure strategy profiles. A pure Nash equilibrium
is a pure strategy profile (s, t) such that u(s , t) u(s, t) for all s S and v(s, t )
v(s, t) for all t T .
To define the mixed extension we need notational conventions for probability
measures on finite sets. For each k = 0, 1, 2, . . . let
k1 = { Rk+ : 1 + + k = 1 }
be the k 1 dimensional simplex. We will typically think of this as the set of
probability measures on a set with k elements indexed by the integers 1, . . . , k. In
particular, let S = m1 and T = n1 ; elements of these sets are called mixed
strategies for agents 1 and 2 respectively. Abusing notation, we will frequently
identify pure strategies si S and tj T with the mixed strategies in S and T
that assign all probability to i and j.
An element of S T is called a mixed strategy profile. We let u and v
also denote the bilinear extensions of the given payoff functions to S T , so the
expected payoffs resulting from a mixed strategy profile (, ) S T are
u(, ) =
m X
n
X
u(si , tj )i j
and
i=1 j=1
v(, ) =
m X
n
X
v(si , tj )i j
i=1 j=1
and v(, ) = T B,
1 ( ) = argmax A
S
and
2 () = argmax T B .
T
38
3 3 4
4 5
A= 4 3 3
and
B= 4 2
3 4 3
5 4
m = n = 3, with
2
5 .
2
s2
t2
t3
T
s3
t1
s2
t2
s1
s3
s1
t1
Figure 3.1
t3
39
A similar procedure can be used to find Nash equilibria in which each agent
mixes over two pure strategies. If we consider s1 and s2 , we see that there are two
mixtures that allow agent 2 to mix over two pure strategies, and we will need to
consider both of them, so things are a bit more complicated than they were for pure
strategies because the process branches. Suppose that agent 1 mixes over s1 and
s2 in the proportion that makes t1 and t2 best responses. Agent 2 has a mixture of
t1 and t2 that makes s2 and s3 best responses. There is a mixture of s2 and s3 that
makes t1 and t3 best responses, and a certain mixture of t1 and t3 makes s1 and
s2 best responses. The only hope for continuing this path in a way that might lead
to a Nash equilibrium is to now consider the mixture of s1 and s2 that makes t1
and t3 best responses, and indeed, ( , ) is a Nash equilibrium.
We havent yet considered the possibility that agent 1 might mix over s1 and s3 ,
nor have we examined what might happen if agent 2 mixes over t2 and t3 . There is
a mixture of s1 and s3 that allow agent 2 to mix over t1 and t2 , which is a possibility
we have already considered and there is a mixture of t2 and t3 that allows agent
1 to mix over s1 and s3 , which we also analyzed above. Therefore there are no
additional Nash equilibria in which both agents mix over two pure strategies.
Could there be a Nash equilibrium in which one of the agents mixes over all
three pure strategies? Agent 2 does have one mixed strategy that allows agent 1 to
mix freely, but this mixed strategy assigns positive probability to all pure strategies
(such a mixed strategy is said to be totally mixed) so it is not a best response
to any of agent 1s mixed strategies, and we can conclude that there is no Nash
equilibrium of this sort. Thus ( , ) is the only Nash equilibrium.
This sort of analysis quickly becomes extremely tedious as the game becomes
larger. In addition, the fact that we are able to find all Nash equilibria in this way
does not prove that there is always something to find.
Before continuing we reformulate Nash equilibrium using a simple principle with
numerous repercussions, namely that a mixed strategy maximizes expected utility if
and only if it assigns all probability to pure strategies that maximize expected utility.
To understand this formally it suffices to note that agent 1s problem is to maximize
T
ui (, ) = A =
m
X
i=1
n
X
j=1
Pm
i=1
aij j
Lemma 3.1.1. A mixed strategy profile (, ) is a Nash equilibrium if and only if:
Pn
Pn
(a) for each i = 1, . . . , m, either i = 0 or
j=1 ai j j for all
j=1 aij j
i = 1, . . . , m;
P
Pm
ij
i
i=1
i=1 bij i for all j =
1, . . . , n.
For each m + n conditions there are two possibilities, so there are 2m+n cases. For
each of these cases the intuition derived from counting equations and unknowns
40
suggests that the set of solutions of the conditions given in Lemma 3.1.1 will typically be zero dimensional, which is to say that it is a finite set of points. Thus we
expect that the set of Nash equilibria will typically be finite.
The Lemke-Howson algorithm is based on the hope that if we relax one of the
conditions above, say the one saying that either 1 = 0 or agent 1s first pure
strategy is a best response, then we may expect that the resulting set will be one
dimensional. Specifically, we let M be the set of pairs (, ) S T satisfying:
P
P
(a) for each i = 2, . . . , m, either i = 0 or nj=1 aij j nj=1 ai j j for all i =
1, . . . , m;
Pm
P
s2
t2
t3
Db
s3
B
t1
t2
s1
s2
A
b
s3
b
b
t1
s1
E
b
t3
Figure 3.2
41
best response at A, so there is the possibility of holding A fixed and moving away
from t2 along the edge of T between t1 and t2 . We cant continue in this way past
B because s3 would no longer be a best response. However, at B both s2 and
s3 are best responses, so the conditions defining M place no constraints on agent
1s mixed strategy. Therefore we can move away from (A, B) by holding B fixed
and moving into the interior of S in a way that obeys the constraints on agent 2s
mixed strategy, which are that t1 and t2 are best responses. This edge bumps into
the boundary of S at C. Since the probability of s3 is now zero, we are no longer
required to have it be a best response, so we can continue from B along the edge of
T until we arrive at t1 . Since the probability of t2 is now zero, we can move away
from C along the edge between s1 and s2 until we arrive at D. Since t3 is now a
best response, we can move away from t1 along the edge between t1 and t3 until we
arrive at E. As we saw above, (D, E) = ( , ) is a Nash equilibrium.
We now explain how this works in general. If Y is a proper subset of {1, . . . , m}
and D is a nonempty subset of {1, . . . , n}, let
SY (D) = { S : i = 0 for all i Y and D argmax
j=1,...,n
X
i
bij i }
be the set of mixed strategies for agent 1 that assign zero probability to every pure
strategy in Y and make every pure strategy in D a best response. Evidently SY (D)
is a polytope.
It is now time to say what typically means. The matrix B is said to be in
Lemke-Howson general position if, for all Y and D, SY (D) is either empty or
(m |D| |Y |)-dimensional. That is, SY (D) has the dimensions one would expect
by counting equations and unknowns. In particular, if m < |D| + |Y |, then SY (D)
is certainly empty.
Similarly, if Z is a proper subset of {1, . . . , n} and C is a nonempty subset of
{1, . . . , m}, let
TZ (C) = { T : j = 0 for all j Z and C argmax
i=1,...,m
X
j
aij j }.
The matrix A is said to be in Lemke-Howson general position if, for all Z and C,
TZ (C) is either empty or (n |C| |Z|)-dimensional. Through the remainder of
this section we assume that A and B are in Lemke-Howson general position.
The set of Nash equilibria is the union of the cartesian products SY (D) TZ (C)
over all quadruples (Y, D, Z, C) with Y C = {1, . . . , m} and Z D = {1, . . . , n}.
The general position assumption implies that if such a product is nonempty, then
|Y | + |C| = m and |Z| + |D| = n, so that Y and C are disjoint, as are Z and D,
and SY (D) TZ (C) is zero dimensional, i.e., a singleton. Thus the general position
assumption implies that there are finitely many equilibria.
In addition, we now have
[
M=
SY (D) TZ (C)
()
where the union is over all quadruples (Y, D, Z, C) such that:
42
43
is the starting point of the algorithm, in which case it is not an endpoint of any
edge quadruple.
Now suppose that (Y, D, Z, C) is not a Nash equilibrium. Since SY (D) and
TZ (C) are 0-dimensional, |D| + |Y | = m and |C| + |Z| = n, so, in view of (e), one
of the two intersections Y C and Z D is a singleton while the other is empty.
First suppose that Z D = {j}. Then (Y, D, Z \ {j}, C) and (Y, D \ {j}, Z, C)
are the only quadruples that might be edge quadruples that have (Y, D, Z, C) as an
endpoint, and in fact both are: again (a)-(d) hold obviously (except that one must
note that |D| 2 because |Z D| = n, |Z| < n, and |Z D| = 1) and SY (D \ {j})
and TZ\{j} (C) are both nonempty because SY (D) and TZ (C) are nonempty subsets.
On the other hand, if Y C = {i}, then (Y \ {i}, D, Z, C) and (Y, D, Z, C \ {i})
are the only quadruples that might be edge quadruples that have (Y, D, Z, C) as an
endpoint. By the logic above, (Y \ {i}, D, Z, C) certainly is, and (Y, D, Z, C \ {i})
is if C 6= {i}, and not otherwise. When C = {i} we have Y C = {2, . . . , m}
and Y C = {i} = C, so Y = {2, . . . , m}, which is to say that (Y, D, Z, C) is the
starting point of the algorithm. In sum, if (Y, D, Z, C) is not a Nash equilibrium, it
is an endpoint of precisely two edge quadruples except when it is the starting point
of the algorithm, in which case is an endpoint of precisely one edge quadruple.
Taken together, these observations verify (i)-(iii), and complete the formal verification of the main properties of the Lemke-Howson algorithm. Two aspects of
the procedure are worth noting. First, when SY (D) TZ (C) is a vertex that
is an endpoint of two edges, the two edges are either SY \{i} (D) TZ (C) and
SY (D) TZ (C \ {i}) for some i or SY (D) TZ\{j} (C) and SY (D \ {j}) TZ (C) for
some j. In both cases one of the edges is the cartesian product of a line segment
in S and a point in T while the other is the cartesian product of a point in S and
a line segment in T . Geometrically, the algorithm alternates between motion in S
and motion in T .
Second, although our discussion has singled out the first pure strategy of agent
1, this was arbitrary, and any pure strategy of either player could be designated for
this role. It is quite possible that different choices will lead to different equilibria.
In addition, although the algorithm was described in terms of starting at this pure
strategy and its best response, the path following procedure can be started at any
endpoint of a path in M. In particular, having computed a Nash equilibrium using
one designated pure strategy, we can then switch to a different designated pure
strategy and follow the path, for the new designated pure strategy, going away
from the equilibrium. This path may go to the starting point of the algorithm
for the new designated pure strategy, but it is also quite possible that it leads
to a Nash equilibrium that cannot be reached directly by the algorithm using any
designated pure strategy. Equilibria that can be reached by repeated applications of
this maneuver are said to be accessible. A famous example due to Robert Wilson
(reported in Shapley (1974))) shows that there can be inaccessible equilibria even
in games with a surprisingly small number of pure strategies.
44
3.2
45
2
2
S
t3
s3
t1
s1
3
s2
t2
1
1
Figure 3.3
We now transport the Lemke-Howson algorithm to this framework. Let M be
the set of (, ) S T such that, when we set s = em A and t = en B T ,
we have
(a) for each i = 2, . . . , m, either i = 0 or si = 0;
(b) for each j = 1, . . . , n, either j = 0 or tj = 0.
For our running example we can follow a path in M from (0, 0) to the image of
the Nash equilibrium, as shown in Figure 3.4. This path has a couple more edges
than the one in Figure 3.2, but there is the advantage of starting at (0, 0), which is
a bit more canonical.
If we set
0 A
,
q = e , y = (, ), and x = (s, t),
= m + n,
C=
BT 0
the system above is equivalent to
Cy + x = q
hx, yi = 0
x, y 0 R .
()
46
We will assume that all the components of q are positive, that all the entries of C
are nonnegative, and that each row of C has at least one positive entry, so that P
is bounded and thus a polytope. In general a d-dimensional polytope is said to be
simple if each of its vertices is in exactly d facets. The condition that generalizes
the general position assumption on A and B is that P is simple.
Let the projection of P onto the second copy of R be
Q = { y R : y 0 and Cy q }.
0 A
and q = e , then Q = S T , and each edge of Q is either the
If C =
BT 0
cartesian product of a vertex of S and an edge of T or the cartesian product of
an edge of S and a vertex of T .
2
2
t3
s3
b
t1
s1
s2
b
t2
1
1
Figure 3.4
Our problem is to find a (x, y) P such that x 6= 0 satisfying the complementary slackness condition hx, yi = 0. The algorithm follows the path starting at
(x, y) = (q, 0) in
M = { (x, y) P : x2 y2 + + x y = 0 }.
The equation x2 y2 + +x y = 0 encodes the condition that for each j = 2, . . . , , either xj = 0 or yj = 0. Suppose we are at a vertex (x, y) of P satisfying this condition,
47
1
1
ci2
ci
qi xi y2 y .
ci1
ci1
ci1
ci1
Replacing the first equation above with this, and substituting it into the other
equations, gives
c11 ci2
c11 ci
c11 c11
xi c12
y2 c1
y ,
x1 = q1
qi
ci1
ci1
ci1
ci1
..
..
..
..
..
.
.
.
.
.
1
ci2
1
qi
xi
y2
ci1
ci1
ci1
..
..
..
.
.
.
c1 c1
c1 ci2
= q
xi c2
y2 c
qi
ci1
ci1
ci1
y1 =
..
.
x
ci
y ,
ci1
..
.
c1 ci
y .
ci1
This is not exactly a thing of beauty, but it evidently has the same form as what
we started with. The data of the algorithm consists of a tableau [q , C ], a list
describing how the rows and the last columns of the tableau correspond to the
original variables of the problem, and the variable that vanished when we arrived
at the corresponding vertex. If this variable is either x1 or y1 we are done. Otherwise the data is updated by letting the variable that is complementary to this one
48
increase, finding the next variable that will vanish when we do so, then updating
the list and the tableau appropriately. This process is called pivoting.
We can now describe how the algorithm works in the degenerate case when P
is not necessarily simple. From a conceptual point of view, our method of handling
degenerate problems is to deform them slightly, so that they become nondegenerate,
but in the end we will have only a combinatoric rule for choosing the next pivot
variable. Let L = { (x, y) R R : Cy + x = q }, let 1 , . . . , , 1 , . . . , be
distinct positive integers, and for > 0 let
P = { (x, y) L : xi i and yi i for all i = 1, . . . , }.
If (x, y) is a vertex of P , then there are variables, which we will describe as
free variables, whose corresponding equations xi = i and yi = i determine
(x, y) as the unique member of L satisfying them. At the point in L where these
equations are satisfied, the other variables can be written as linear combinations of
the free variables, and thus as polynomial functions of . Because the i and i
are all different, there are only finitely many values of such that any of the other
variables vanish at this vertex. Because there are finitely many -element subsets
of the 2 variables, it follows that P is simple for all but finitely many values of .
In particular, for all in some interval (0, ) the combinatoric structure of P will
be independent of . In addition, we do not actually need to work in P because the
pivoting procedure, applied to the polytope P for such , will follow a well defined
path that can be described in terms of a combinatoric procedure for choosing the
next pivot variable.
To see what we mean be this consider the problem of finding which xi first goes
below i as we go out the line y1 1 , y2 = 2 , . . . , y = . This is
basically a process of elimination. If ci1 0, then increasing y1 never leads to a
violation of the ith constraint, so we can begin by eliminating all those i for which
ci1 is not positive. Among the remaining i, the problem is to find the i for which
1
ci2
ci
1
qi + i + 2 + +
ci1
ci1
ci1
ci1
is smallest for small > 0. The next step is to eliminate all i for which qi /c1i is not
minimal. For each i that remains the expression
1 i ci2 2
ci
+ + +
ci1
ci1
ci1
has a dominant term, namely the term, among those with nonzero coefficients,
whose exponent is smallest. The dominant terms are ordered according to their
values for small > 0:
(a) terms with positive coefficients are greater than terms with negative coefficients;
(b) among terms with positive coefficients, those with smaller exponents are
greater than terms with larger exponents, and if two terms have equal exponents they are ordered according to the coefficients;
49
(c) among terms with negative coefficients, those with larger exponents are greater
than terms with smaller exponents, and if two terms have equal exponents
they are ordered according to the coefficients.
We now eliminate all i for which the dominant term is not minimal. All remaining
i have the same dominant term, and we continue by subtracting off this term and
comparing the resulting expressions in a similar manner, repeating until only one i
remains. This process does necessarily continue until only one i remains, because
if other terms of the expressions above fail to distinguish between two possibilities,
eventually there will be a comparison involving the terms i /ci1 , and the exponents
1 , . . . , , 1 , . . . , are distinct.
Lets review the situation. We have given an algorithm that finds a solution
of the linear complementarity problem () that is different from (q, 0). The assumptions that insure that the algorithm works are that q 0 and that P is
a polytope. In particular, these assumptions are satisfied when the linear complementarity problem is derived from a two person game with positive payoffs, in which
case any solution other than (q, 0) corresponds to a Nash equilibrium. Therefore
any two person game with positive payoffs has a Nash equilibrium, but since the
equilibrium conditions are unaffected by adding a constant to a players payoffs, in
fact we have now shown that any two person game has a Nash equilibrium.
There are additional issues that arise in connection with implementing the algorithm, since computers cannot do exact arithmetic on arbitrary real numbers.
One possibility is to require that the entries of q and C lie in a set of numbers
for which exact arithmetic is possibleusually the rationals, but there are other
possibilities, at least theoretically. Alternatively, one may work with floating point
numbers, which is more practical, but also more demanding because there are issues
associated with round-off error, and in particular its accumulation as the number of
pivots increases. The sort of pivoting we have studied here also underlies the simplex algorithm for linear programming, and the same sorts of ideas are applied to
resolve degeneracy. Numerical analysis for linear programming has a huge amount
of theory, much of which is applicable to the Lemke-Howson algorithm, but it is far
beyond our scope.
3.3
50
Of course this observation does not prove anything, but it does point in a useful
direction. Let x1 , . . . , xn , y1 , . . . , yn X be given. We can define a finite two person
game with n n payoff matrices A = (aij ) and B = (bij ) by setting
(
0, i 6= j,
aij = kxi yj k2 and bij =
1, i = j.
Let (, ) n1 n1 be a mixed strategy profile. Clearly is a best response to
if and only if it assigns all probability to the strategies that are assigned maximum
probability by , which is to say that j > 0 implies that j i for all i.
Understanding
when is a best response to requires a brief calculation. Let
P
z = nj=1 j yj . For each i we have
X
j
aij j =
X
j
j kxi yj k2 =
j xi yj , xi yj
j
=
j xi , xi + 2
j xi , yj
j yj , yj
j
= xi , xi + 2 xi , z hz, zi + C = kxi zk2 + C
P
where C = kzk2 nj=1 j kyj k2 is a quantity that does not depend on i. Therefore
is a best response to if and only if it assigns all probability to those i with xi as
close to z as possible. If y1 F (x1 ), . . . , yn F (xn ), then there is a sense in which
a Nash equilibrium may be regarded as a point that is approximately fixed.
We are going to make this precise, thereby proving Kakutanis fixed point theorem. Assume now that F is upper semicontinuous with convex values. Define
sequences x1 , x2 , . . . and y1 , y2 , . . . inductively as follows. Choose x1 arbitrarily, and
let y1 be an element of F (x1 ). Supposing that x1 , . . . , xn and y1 , . . . , yn , have already been determined, let ( n , n ) be a Nash equilibrium of the two person game
with payoff matrices An = (anij ) and B nP= (bnij ) where anij = kxi yj k2 and bnij is
1 if i = j and 0 otherwise. Let xn+1 = j j yj , and choose yn+1 F (yn+1).
Let x be an accumulation point of the sequence {xn }. To show that x is a
fixed point of F it suffices to show that it is an element of the closure of any convex
neighborhood V of F (x ). Choose
0 such that F (x) V for all x U (x ).
P >
Consider an n such that xn+1 = j jn yj U/3 (x ) and at least one of x1 , . . . , xn
is also in this ball. Then the points in x1 , . . . , xn that are closest to xn+1 are in
U2/3 (xn+1 ) U (x ), so xn+1 is a convex combination of points in V , and is
therefore in V . Therefore x is in the closure of the set of xn that lie in V , and thus
in the closure of V .
In addition to proving the Kakutani fixed point theorem, we have accumulated
all the components of an algorithm for computing approximately fixed points of
a continuous function f : X X. Specifically, for any error tolerance > 0 we
compute the sequences x1 , x2 , . . . and y1 , y2 , . . . with f in place of F , halting when
kxn+1 f (xn+1 )k < . The argument above shows that this is, in fact, an algorithm,
in the sense that it is guaranteed to halt eventually. This algorithm is quite new.
Code implementing it exists, and the initial impression is that it performs quite
well. But it has not been extensively tested.
51
There is one more idea that may have some algorithmic interest. As before, we
consider points x1 , . . . , xn , y1 , . . . , yn Rd . Define a correspondence : Rd Rd
by letting (z) be the convex hull of { yj : j argmini kz xi k } when z PJ .
(Evidently this construction is closely related to the Voronoi diagram determined by
x1 , . . . , xn . Recall that this is the polyhedral decomposition of Rd whose nonempty
polyhedra are the sets PJ = { z V : J argmini kz xi k } where =
6 J
{1, . . . , n}.) Clearly is upper semicontinuous and convex valued.
Suppose that
P z is a fixed point of this correspondence. Then z is a convex
combination
/ argmini kz xi k. Let J = { j : yj >
j j yj with yj = 0 if j
0 }. If i = 1/|J| when i J and i = 0 when i
/ J, then (, ) is a Nash
equilibrium of the game derived from xP
,
.
.
.
,
x
,
y
,
.
.
.
, yn . Conversely, if (, ) is
1
n 1
a Nash equilibrium of this game, then jJ j yj is a fixed point of . In a sense,
the algorithm described above approximates the given correspondence F with a
correspondence of a particularly simple type.
We may project the path of the Lemke-Howson algorithm, in its application to
the game derived from x1 , . . . , xn , y1, . . . , yn , into this setting. Define 1 : Rd Rd
by letting 1 (z) be the convex hull of { yi : i {1}argmini kzxi k }. Suppose that
(, ) is an element of the set M defined in Section 3.1, so that all the conditions
of Nash equilibrium are satisfied except that it may be the case that 1 > 0Peven if
the first pure strategy is not optimal. Let J = { j : j > 0 }, and let z = j j yj .
Then J { i : i > 0 } {1} argminj kz xj k, so z 1 (z). Conversely,
suppose
P
z is a fixed point of 1 , and let J = argminj kz xj k. Then z = j j yj for some
n1 with j = 0 for all j
/ {1} J. If we let be the element of n1 such
that i = 1/|{1} J| if i J and i = 0 if i
/ {1} J, then (, ) M.
If n is large one might guess that there is a sense in which operating in Rd might
be less burdensome than working in n1 n1 , but it seems to be difficult to
devise algorithms that take concrete advantage of this. Nonetheless this setup does
give a picture of what the Lemke-Howson algorithm is doing that has interesting
implications. For example, if there is no point in Rd that is equidistant from more
than d + 1 points, then there is no point (, ) M with i > 0 for more than
d + 2 indices. This gives a useful upper bound on the number of pivots of the
Lemke-Howson algorithm.
3.4
Sperners Lemma
52
3
b
3
b
1
b
1
3
1
b
b
b
2
b
Figure 3.5
Theorem 3.4.1 (Sperners Lemma). If is a Sperner labelling, then the number
of completely labelled simplices is odd.
Before proving this, lets see why its important:
Proof of Browers Theorem. Let f : d1 d1 be a continuous function. Proposition 2.5.2 implies that there is a sequence P1 , P2 , . . . of triangulations whose meshes
converge to zero. For each r = 1, 2, . . . let V r be the set of vertices of Pr . If any
of the elements of V r is a fixed point we are done, and otherwise we can define
r : V r {0, . . . , d} by letting r (v) be the smallest index i such that vi > fi (v).
Evidently r is a Sperner labelling, so there is a completely labelled simplex with
vertices v1r , . . . , vdr where r (vir ) = i. Passing to a subsequence, we may assume that
the sequences vi1 , vi2 , . . . have a common limit x. For each i we have
fi (x) = lim fi (v r ) lim vir = xi ,
and
fi (x) = 1 =
xi , so f (x) = x.
We will give two proofs of Sperners lemma. The first of these uses facts about
volume, and in this sense is less elementary than the second (which is given in the
next section) but it quickly gives both an intuition for why the result is true and
an important refinement.
We fix an affine isometry1 A : H d1 Rd1 such that
D = det A(e2 ) A(e1 ), . . . , A(ed ) A(e1 ) > 0.
1
If (X, dX ) and (Y, dY ) are metric spaces, a function : X Y is an isometry if
dY ((x), (x )) = dX (x, x ) for all x, x X.
53
1
det A(v2 (t)) A(v1 (t)), . . . , A(vd (t)) A(v1 (t)) .
d!
For 0 t 1 let (t) be the convex hull of v1 (t), . . . , vd (t). Then p (t) is the
volume of (t) when t is small.
We have
1
p (1) = det A(e(v2 ) ) A(e(v1 ) ), . . . , A(e(vd ) ) A(e(v1 ) ) .
d!
Elementary properties of the determinant imply that each p and p are polynomial
functions. For sufficiently small t the simplices (t) are the (d 1)-dimensional
simplices of a triangulation of d1 .3 Therefore p(t) is d!1 D for small t. Since p is a
2
Actually, it is straightforward if you know integration, but Gauss regarded this as too heavy
a tool, expressing a wish for a more elementary theory of the volume of polytopes. The third of
Hilberts famous problems asks whether it is possible, for any two polytopes of equal volume, to
triangulate the first in such a way that the pieces can be reassembled to give the second. This
was resolved negatively by Hilberts student Max Dehn within a year of Hilberts lecture laying
out the problems, and it remains the case today that there is no truly elementary theory of the
volumes of polytopes. In line with this, our discussion presumes basic facts about d-dimensional
measure of polytopes in Rd that are very well understood by people with no formal mathematical
training, but which cannot be justified formally without appealing to relatively advanced theories
of measure and integration.
3
This is visually obvious, and a formal proof would be tedious, so we provide only a sketch.
Suppose that for each v V we have a path connected neighborhood Uv of v in the interior of the
smallest face of d1 containing v, and this system of neighborhoods satisfies the condition that
for any simplex in P, say with vertices v1 , . . . , vk , if v1 Uv1 , . . . , vk Uvk , then v1 , . . . , vk are
affinely independent. We claim that a simplicial complex obtained by replacing each v with some
element of Uv is a triangulation of d1 ; note that this can be proved by moving one vertex at a
time along a path. Finally observe that because is a Sperner labelling, for each v and 0 t < 1,
v(t) is contained in the interior of the smallest face of d1 containing v.
54
1
D.
d!
3.5
55
3
b
3
b
1
b
bb
bb
1
3
1
b
bb
bb
bb
bb
bb
bb
2
b
bb
bb
bb
bb
bb
Figure 3.6
Of course for each k = 0, . . . , d 2 the set of simplices in P that lie in k
56
and
E = E 0 F 1 E 1 E d2 F d1 E d1 .
57
V, we can compute the simplices of its neighbors in and the labels of the vertices
of these simplices. If we remember which of these neighbors we were at prior to
arriving at the current element of V, then the next step in the algorithm is to go to
the other neighbor. Such a step along the path of the algorithm is called a pivot.
3
b
3
b
3
b
1
b
bb
bb
1
3
1
b
bb
bb
bb
bb
bb
bb
2
bb
bbb
bb b
bbb
bbb
bb b
Figure 3.7
At this point we remark on a few aspects of the Scarf algorithm, and later
we will compare it with various alternatives. The first point is that it necessarily
moves through d1 rather slowly. Consider a k-almost completely labelled simplex
. Each pivot of the algorithm drops one of the vertices of the current simplex,
possibly adding a new vertex, or possibly dropping down to a lower dimensional
face. Therefore a minimum of k pivots are required before one can possibly arrive
at a simplex that has no vertex in common with . If the grid is fine, the algorithm
will certainly require many pivots to arrive a fixed point far from the algorithms
starting point.
This suggests the following strategy. We first apply the Scarf algorithm to a
coarse given triangulation of d1 , thereby arriving at a completely labelled simplex
that is hopefully a rough approximation of a fixed point. We then subdivide the
given triangulation of d1 , using barycentric subdivision or some other method.
If we could somehow restart the algorithm in the fine triangulation, near the
completely labelled simplex in the coarse triangulation, it might typically be the
case that the algorithm did not have to go very far to find a completely labelled
simplex in the fine triangulation. Restart methods do exist (see, e.g., Merrill (1972),
Kuhn and MacKinnon (1975), and van der Laan and Talman (1979)) but it remains
the case that the Scarf algorithm has not proved to be very useful in practice,
perhaps due in part to its difficulties with high dimensional problems.
There is one more feature of the Scarf algorithm that is worth mentioning. In
our description of the algorithm the ordering of the vertices plays an explicit role,
and can easily make a difference to the outcome. If one wishes to find more than
one completely labelled simplex, or perhaps as many as possible, or perhaps even all
of them, there is the following strategy. Having followed the algorithm for the given
ordering of the indices to its terminus, now proceed from that completely labelled
simplex in the graph associated with some different ordering. This might lead
58
back to the starting point of the algorithm in , but it is also quite possible that
it might lead to some completely labelled simplex that cannot be reached directly
by the algorithm under any ordering of the indices. A completely labelled simplex
is accessible if it is reachable by the algorithm in this more general sense: there
is path going to from the starting point of the algorithm for some ordering of the
indices, along a path that is a union of maximal paths of the various graphs for
the various orderings of the indices.
3.6
Homotopy
59
3.7
Remarks on Computation
We have now seen three algorithms for computing points that are approximately
fixed. How good are these, practically and theoretically? The first algorithm we saw,
in Section 3.3, is new. It is simple, and can be applied to a wide variety of settings.
Code now exists, but there has been little testing or practical experience. The Scarf
algorithm has not lived up to the hopes it raised when it was first developed, and is
not used in practical computation. Homotopy methods are restricted to problems
that are smooth. As we mentioned above, within this domain they have an extensive
track record with considerable success.
More generally, what can we reasonably hope for from an algorithm that computes points that are approximately fixed, and what sort of theoretical concepts can
we bring to bear on these issues? These question has been the focus of important
recent advances in theoretical computer science, and in this section we give a brief
description of these developments. The discussion presumes little in the way of
prior background in computer science, and is quite superficiala full exposition of
this material is far beyond our scope. Interested readers can learn much more from
the cited references, and from textbooks such as Papadimitriou (1994a) and Arora
and Boaz (2007).
Theoretical analyses of algorithms must begin with a formal model of computation. The standard model is the Turing machine, which consists of a processor with
finitely many states connected by an input-output device to a unbounded one dimensional storage medium that records data in cells, on each of which one can write
an element of a finite alphabet that includes a distinguished character blank. At
the beginning of the computation the processor is in a particular state, the storage
medium has a finitely many cells that are not blank, and the input-output device
is positioned at a particular cell in storage. In each step of the computation the
character at the input-output devices location is read. The Turing machine is essentially defined by functions that take state-datum pairs as their arguments and
60
compute:
the next state of the processor,
a bit that will be written at the current location of the input-output device
(overwriting the bit that was just read) and
a motion (forward, back, stay put) of the input-output device.
The computation ends when it reaches a particular state of the machine called
Halt. Once that happens, the data in the storage device is regarded as the
output of the computation.
As you might imagine, an analysis based on a concrete and detailed description
of the operation of a Turing machince can be quite tedious. Fortunately, it is
rarely necessary. Historically, other models of computation were proposed, but were
subsequently found to be equivalent to the Turing model, and the Church-Turing
thesis is the hypothesis that all reasonable models of computation are equivalent,
in the sense that they all yield the same notion of what it means for something to be
computable. This is a metamathematical assertion: it can never be proved, and a
refutation would not be logical, but would instead be primarily a social phenomenon,
consisting of researchers shifting their focus to some inequivalent model.
Once we have the notion of a Turing machine, we can define an algorithm to
be a Turing machine that eventually halts, for any input state of the storage device.
A subtle distinction is possible here: a Turing machine that always halts is not
necessarily the same thing as a Turing machine that can be proved to halt, regardless
of the input. In fact one of the most important early theorems of computer science
is that there is no algorithm that has, as input, a description of a Turing machine
and a particular input, and decides whether the Turing machine with that input will
eventually halt. As a practical matter, one almost always works with algorithms
that can easily be proved to be such, in the sense that it is obvious that they
eventually halt.
A computational problem is a rule that associates a nonempty set of outputs
with each input, where the set of possible inputs and outputs is the set of pairs
consisting of a position of the input-output device and a state of the storage medium
in which there are finitely many nonblank cells. (Almost always the inputs of
interest are formatted in some way, and this definition implicitly makes checking the
validity of the input part of the problem.) A computational problem is computable
if there is an algorithm that passes from each input to one of the acceptable outputs.
The distinction between computational problems that are computable and those
that are not is fundamental, with many interesting and important aspects, but in
our discussion here we will focus exclusively on problems that are known to be
computable.
For us the most important distinctions is between those computable computational problems that are easy and those that are hard, where the definitions
of these terms remain to be specified. In order to be theoretically useful, the easiness/hardness distinction should not depend on the architecture of a particular
machine or the technology of a particular era. In addition, it should be robust, at
least in the sense that a composition of two easy computational problems, where
61
the output of the first is the input of the second, should also be easy, and possibly in other senses as well. For these reasons, looking at the running time of an
algorithm on a particular input is not very useful. Instead, it is more informative
to think about how the resources (time and memory) consumed by a computation
increase as the size of the input grows. In theoretical computer science, the most
useful distinction is between algorithms whose worst case running time is bounded
by a polynomial function of the size of the output, and algorithms that do not
have this property. The class of computational problems that have polynomial time
algorithms is denoted by P. If the set of possible inputs of a computational problem is finite, then the problem is trivially in P, and in fact we will only consider
computational problems with infinite sets of inputs.
There are many kinds of computational problems, e.g., sorting, function evaluation, optimization, etc. For us the most important types are decision problems ,
which require a yes or no answer to a well posed question, and search problems,
which require an instance of some sort of object or a verification that no such object exists. An important example of a decision problem is Clique: given a simple
undirected graph G and an integer k, determine whether G has a clique with k
nodes, where a clique is a collection of vertices such that G has an edge between
any two of them. An example of a search problem is to actually find such a clique
or to certify that no such clique exists.
There is a particularly important class of decision problems called NP, which
stands for nondeterministic polynomial time. Originally NP was thought of as
the class of decision problems for which a Turing machine that chose its next state
randomly has a positive probability of showing that the answer is Yes when this
is the case. For example, if a graph has a k-clique, an algorithm that simply guesses
which elements constitute the clique has a positive probability of stumbling onto
some k-clique. The more modern way of thinking about NP is that it is the class of
decision problems for which a Yes answer has a certificate or witness that can
be verified in polynomial time. In the case of Clique an actual k-clique is such a
witness. Factorization of integers is another algorithmic issue which easily generates
decision problemsfor example, does a given number have a prime factor whose
first digit is 3?that are in NP because a prime factorization is a witness for them.
(One of the historic recent advances in mathematics is the discovery of a polynomial
time algorithm for testing whether a number is prime. Thus it is possible to verify
the primality of the elements of a factorization in polynomial time.)
An even larger computational class is EXP, which is the class of computational
problems that have algorithms with running times that are bounded above by a
function of the form exp(p(s)), where s is the size of the problem and p is a polynomial function. Instead of using time to define a computational class, we can
also use space, i.e., memory; PSPACE is the class of computational problems that
have algorithms that use an amount of memory that is bounded by a polynomial
function of the size of the input. The sizes of the certificates for a problem in
NP are necessarily bounded by some polynomial function of the size of the input,
and the problem can be solved by trying all possible certificates not exceeding this
bound, so any problem in NP is also in PSPACE. In turn, the number of processor
state-memory state pairs during the run of a program using polynomially bounded
62
63
naturally; Clique is one of them. One of the most famous problems in contemporary mathematics is to determine whether NP is contained in P. This question
boils down to deciding whether Clique (or any other NP-complete problem) has
a polynomial time algorithm. This is thought to be highly unlikely, both because a
lot of effort has gone into designing algorithms for these problems, and because the
existence of such an algorithm would have remarkable consequences. It should be
mentioned that this problem is, to some extent at least, an emblematic representative of numerous open questions in computer science that have a similar character.
In fact, one of the implicit conventions of the discipline is to regard a computational
problem as hard if, after some considerable effort, people havent been able to figure
out whether it is hard or easy.
For any decision problem in NP there is an associated search problem, namely
to find a witness for an affirmative answer or verify that the answer is negative.
For Clique this means not only showing that a clique of size k exists, but actually
producing one. The class of search problems associated with decision problems is
called FNP. (The F stands for function.) For Clique the search problem is
not much harder than the decision problem, in the following sense: if we had a
polynomial time algorithm for the decision problem, we could apply it to the graph
with various vertices removed, repeatedly narrowing the focus until we found the
desired clique, thereby solving the search problem is polynomial time.
However, there is a particular class of problems for which the search problem
is potentially quite hard, even though the decision problem is trivial because the
answer is known to be yes. This class of search problems is called TFNP. (The
T stands for total.) There are some trivial decision problems that give rise
to quite famous problems in this class:
Does a integer have a prime factorization? Testing primality can now be
done in polynomial time, but there is still no polynomial time algorithm for
factoring.
Given a set of positive integers {a1 , . . . , an } with ai < 2n /n for all i, do there
exist two different subsets with the same sum? There are 2n different subsets,
and the sum of any one of them is less than 2n n + 1, so the pigeonhole
principle implies that the answer is certainly yes.
Does a two person game have sets of pure strategies for the agents that are
the supports4 of a Nash equilibrium? Verifying that a pair of sets are the
support of a Nash equilibrium is a computation involving linear algebra and a
small number of inequality verifications that can be performed in polynomial
time.
Problems involving a function defined on some large space must be specified
with a bit more care, because if the function is given by listing its values, then the
problem is easy, relative to the size of the input, because the input is huge. Instead,
one takes the input to be a Turing machine that computes (in polynomial time) the
value of the function at any point in the space.
4
The support of a mixed strategy is the set of pure strategies that are assigned positive
probability.
64
Mainly because the class of problems in NP that always have a positive answer is
defined in terms of a property of the outputs, rather than a property of the inputs
(but also in part because factoring seems so different from the other problems)
experts expect that TFNP does not contain any problems that are complete for
the class. In view of this, trying to study the class as a whole is unlikely to be
very fruitful. Instead, it makes sense to define and study coherent subclasses, and
Papadimitriou (1994b) advocates defining subclasses in terms of the proof that a
solution exists. Thus PPP (polynomial pigeonhole principle) is (roughly) the
class of problems for which existence is guaranteed by the pigeonhole principle, and
PLS (polynomial local search) is (again roughly) the set of problems requesting
a local maximum of a real valued function defined on a graph by a Turing machine.
For us the most important subclass of TFNP is PPAD (polynomial parity
argument directed) which is defined by abstracting certain features of the algorithms we have seen in this chapter. The computational problem EOTL (end of
the line) is defined by a Turing machine that defines a directed graph5 of maximal
degree two in a space that may, without loss of generality, be taken to be the set
{0, 1}k of bit strings of length k, where k is bounded by a polynomial function of
the size of the input. For each v {0, 1}k the Turing machine specifies whether v is
a vertex in the graph. If it is, the Turing machine computes its predecessor, if it has
one, and its successor, if it has one. When it exists, the predecessor of v must be a
vertex, and its successor must be v. Similarly, when v has a successor, it must be a
vertex, and its predecessor must be v. Finally, we require that (0, . . . , 0) is a vertex
that has a successor but no predecessor. The problem is to find another leaf of
the graph, by which we mean either a vertex with a predecessor but no successor,
or a vertex with a successor but no predecessor. Of course the existence of such a
leaf follows from Lemma 2.6.1, generalized in the obvious way to handle directed
graphs. The class of computational problems that have reductions to EOTL is
PPAD (polynomial parity problem directed).
The Lemke-Howson algorithm passes from a two person game to an instance
of EOTL, then solves it by following the path in the graph to its other endpoint.
Similarly, the Scarf algorithm has as input the algorithms for navigating in a triangulation of d1 and generating the labels of the vertices, and if follows a path in
a graph from one endpoint to another. (It would be difficult to describe homotopy
in exactly these terms, but there is an obvious sense in which it has this character.)
5
A directed graph is a pair G = (V, E) where V is a finite set of vertices and E is a finite
set of ordered pairs of elements of V . That is, in a directed graph each edge has a source and a
target.
65
There is a rather subtle point that is worth mentioning here. In our descriptions
of Lemke-Howson, Scarf, and homotopy, we implicitly assumed that the algorithm
used its memory of where it had been to decide which direction to go in the graph,
but the definition of EOTL requires that the graph be directed, which means in
effect that if we begin at any point on the path, we can use local information to decide which of the two directions in the graph constitutes forward motion. It turns
out that each of our three algorithms has this property; a proper explanation of
this would require more information about orientation than we have developed at
this point. The class of problems that can be reduced to the computational problem that has the same features as EOTL, except that the graph is undirected, is
PPA. Despite the close resemblance to PPAD, the theoretical properties of the
two classes differ in important ways.
In a series of rapid developments in 2005 and 2006 (Daskalakis et al. (2006);
Chen and Deng (2006b,a)) it was shown that computing a Nash equilibrium of a
two player game is PPAD-complete, and also that the two dimensional Sperner
problem is PPAD-complete. This means that computing a Nash equilibrium of a
two player game is almost certainly hard, in the sense that there is no polynomial
time algorithm for the problem, because computing general fixed points is almost
certainly hard. Since this breakthrough many other computational problems have
been shown to be PPAD-complete, including finding Walrasian equilibria in seemingly quite simple exchange economies. In various senses the problem does not go
away if we relax the problem, asking for a point that is -approximately fixed for
an that is significantly greater than zero.
The current state of theory presents a contrast between theoretical concepts
that classify even quite simple fixed point problems as intractable, and algorithms
that often produce useful results in a reasonable amount of time. A recent result
presents an even more intense contrast. The computational problem OEOTL has
the same given data as EOTL, but now the goal is to find the other end of the path
beginning at (0, . . . , 0), and not just any second leaf of the graph. Goldberg et al.
(2011) show that OETL is PSPACE-complete, even though the Lemke-Howson
algorithm, the Scarf algorithm, and many specific instances of homotopy procedures
can be recrafted as algorithms for OEOTL.
Recent developments have led to a rich and highly interesting theory explaining
why the problem of finding an approximate fixed point is intractable, in the sense
that there is almost certainly no algorithm that always finds an approximate fixed
point in a small amount of time. What is missing at this point are more tolerant
theoretical concepts that give an account of why the algorithms that exist are as
useful as they are in fact, and how they might be compared with each other, and
with theoretical ideals that have not yet been shown to be far out of reach.
Chapter 4
Topologies on Spaces of Sets
The theories of the degree and the index involve a certain kind of continuity
with respect to the function or correspondence in question, so we need to develop
topologies on spaces of functions and correspondences. The main idea is that one
correspondence is close to another if its graph is close to the graph of the second
correspondence, so we need to have topologies on spaces of subsets of a given space.
In this chapter we study such spaces of sets, and in the next chapter we apply these
results to spaces of functions and correspondences. There are three basic set theoretic operations that are used to construct new functions or correspondences from
given ones, namely restriction to a subdomain, cartesian products, and composition, and our agenda here is to develop continuity results for elementary operations
on sets that will eventually support continuity results for those operations.
To begin with Section 4.1 reviews some basic properties of topological spaces
that hold automatically in the case of metric spaces. In Section 4.2 we define
topologies on spaces of compact and closed subsets of a general topological space.
Section 4.3 presents a nice result due to Vietoris which asserts that for one of these
tolopogies the space of nonempty compact subsets of a compact space is compact.
Economists commonly encounter this in the context of a metric space, in which
case the topology is induced by the Hausdorff distance; Section 4.4 clarifies the
connection. In Section 4.5 we study the continuity properties of basic operations
for these spaces. Our treatment is largely drawn from Michael (1951) which contains
a great deal of additional information about these topologies.
4.1
Topological Terminology
Up to this point the only topological spaces we have encountered have been
subsets of Euclidean spaces. Now it will be possible that X lacks some of the
properties of metric spaces, in part because we may ultimately be interested in some
spaces that are not metrizable, but also in order to clarify the logic underlying our
result.
Throughout this chapter we work with a fixed topological space X. We say that
X is:
(a) a T1 -space if, for each x X, {x} is closed;
66
67
4.2
There will be a number of topologies, and in order to define them we need the
corresponding subbases. For each open U X let:
UU = { K U : K is compact };
UU = UU \ {};
VU = { K X : K is compact and K U 6= };
UU0 = { C U : C is closed };
UU0 = UU0 \ {};
VU0 = { C X : C is closed and C U 6= }.
We now have the following spaces:
K(X)
is the space of compact subsets of X endowed with the topology generated by the subbase { UU : U X is open }.
68
The topologies of H(X) and H0 (X) are both called the Vietoris topology.
H(X) have the subspace topologies induced by the topologies of K (X), K0 (X), and
H0 (X). Of course it is always the case that K(X) and K0 (X) have the subspace
0 (X) respectively.
topologies induced by K(X)
and K
It is easy to see that { UU : U X is open } is a base for K(X) and { UU0 :
U X is open } is a base for K0 (X). Also, for any open U1 , . . . , Uk we have
UU1 . . . UUk = UU1 ...Uk ,
0 (X), and
and similarly for UU , UU0 , and UU0 , so the subbases of K(X),
K(X), K
K0 (X) are actually bases.
4.3
Vietoris Theorem
An interesting fact, which was proved already in Vietoris (1923), and which is
applied from time to time in mathematical economics, is that H(X) is compact
whenever X is compact. We begin the argument with a technical lemma.
Lemma 4.3.1. If X has a subbase such that any cover of X by elements of the
subbase has a finite subcover, then X is compact.
69
Proof. Say that a set is basic if it is a finite intersection of elements of the subbasis.
Any open cover is refined by the collection of basic sets that are subsets of its
elements. If a refinement of an open cover has a finite subcover, then so does the
cover, so it suffices to show that any open cover of X by basic sets has a finite
subcover.
A collection of open covers is a chain if it is completely ordered by inclusion:
for any two covers in the chain, the first is a subset of the second or vice versa. If
each open cover in a chain consists of basic sets, and has no finite subcover, then
the union of the elements of the chain also has these properties (any finite subset
of the union is contained in some member of the chain) so Zorns lemma implies
that if there is one open cover with these properties, then there is a maximal such
cover, say {U : A}.
Suppose, for some A, that U = V1 . . . Vn where V1 , . . . , Vn are in the
subbasis. If, for each i = 1, . . . , n, {U : A} {Vi } has a finite subcover Ci ,
then each Ci \ {Vi } covers X \ Vi , so
(C1 \ {V1 }) . . . (Cn \ {Vn }) {U }
is a finite subcover from {U : A}. Therefore there is at least one i such that
{U : A}{Vi } has no finite subcover, and maximality implies that Vi is already
in the cover. This argument shows that each element U of the cover is contained
in a subbasic set that is also in the cover, so the subbasic sets in {U : A} cover
X, and by hypothesis there must be a finite subcover after all.
Theorem 4.3.2. If X is compact, then H(X) is compact.
Proof. Suppose that { UU : SA} { VV : B} is an open cover of H(X)
by subbasic sets. Let D := X \ V ; since D is closed and X is compact, D is
compact. We may assume that D is nonempty because otherwise X = V1 . . .Vn
for some 1 , . . . , n , in which case H(X) = VV1 . . . VVn . In addition, D must
be contained in some U because otherwise D would not be an element of any UU
or any VV . But then {U } {V : B} has a finite subcover, so, for some
1 , . . . , n , we have
H(X) = UU VV1 . . . VVn .
4.4
Hausdorff Distance
()
70
On the other hand, whenever K U with K compact and U open there is some
> 0 such that U (K) U (otherwise we could take sequences x1 , x2 , . . . in L
and y1 , y2 , . . . in X \ U with d(xi , yi ) 0, then take convergent subsequences) so
{ L : K (L, K) < } UU . Thus:
Lemma 4.4.1. When X is a metric space, the sets of the form { L : K (L, K) < }
constitute a base of the topology of K(X).
The Hausdorff distance between nonempty compact sets K, L X is
H (K, L) := max{K (K, L), K (L, K)}.
This is a metric. Specifically, it is evident that H (K, L) = H (L, K), and that
H (K, L) = 0 if and only if K = L. If M is a third compact set, then
K (K, M) K (K, L) + K (L, M),
from which it follows easily that the Hausdorff distance satisfies the triangle inequality.
There is now an ambiguity in our notation, insofar as U (L) might refer either to
the the union of the -balls around the various points of L or to the set of compact
sets whose Hausdorff distance from L is less than . Unless stated otherwise, we
will always interpret it in the first way, as a set of points and not as a set of sets.
Proposition 4.4.2. The Hausdorff distance induces the Vietoris topology on H(X).
Proof. Fix a nonempty compact K. We will show that any neighborhood of K in
one topology contains a neighborhood in the other topology.
S
First consider some > 0. Choose x1 , . . . , xn K such that K i U/2 (xi ).
If L U/2 (xi ) 6= for all i, then K (L, K) < , so, in view of (),
K UU (K) VU/2 (x1 ) . . . VU/2 (xn ) { L : H (K, L) < }.
We now show that any element of our subbasis for the Vietoris topology contains
{ L : H (K, L) < } for some > 0. If U is an open set containing K, then (as we
argued above) U (K) U for some > 0, so that
K { L : H (L, K) < } { L : K (L, K) < } UU .
If V is open with K V 6= , then we can choose x K V and > 0 small enough
that U (x) V . Then
K { L : H (K, L) < } { L : K (K, L) < } VV .
4.5
71
4.5.1
Continuity of Union
The function taking a pair of sets to their union is as well behaved as one might
hope.
K, H, K
0 , K0 , H0 } the function : (K1 , K2 ) 7
Lemma 4.5.2. For any T {K,
K1 K2 is a continuous function from T (X) T (X) to T (X).
Proof. Applying Lemma 4.5.1, it suffices to show that preimages of subbasic open
K, K
0 , K0 } it suffices to note that
sets are open. For T {K,
1(WU ) = WU WU
U, U 0 , U 0 }. For T {H, H0 } we also need to observe that
for all four W {U,
1(WU ) = (WU H(X)) (H(X) WU )
for both W {V, V 0 }.
4.5.2
Continuity of Intersection
Simple examples show that intersection is not a continuous operation for the
K, K
0 , and K0 .
topologies H and H0 , so the only issues here concern K,
A (X) to K(A)
72
0
For a nonempty closed set A X let KA (X) and KA
(X) be the sets of compact
and closed subsets of X that have nonempty intersection with A. Since the topolo
0 (X),
gies of K(X) and K0 are the subspace topologies inherited from K(X)
and K
last result has the following immediate consequence.
Lemma 4.5.4. The function K 7 K A from KA (X) to K(A) and the function
0
C 7 C A from KA
(X) to K0 (A) are continuous.
Joint continuity of the map (C, D) 7 C D requires an additional hypothesis.
Lemma 4.5.5. If X is a normal space, then : (C, D) 7 C D is a continuous
0 (X) K
0 (X) to K
0 (X). If, in addition, X is a T1 space, then
function from K
: K(X)
K(X)
K(X)
is continuous.
Proof. By Lemma 4.5.1 it suffices to show that, for any open U X, 1 (UU0 ) is
open. For any (C, D) in this set normality implies that there are disjoint open sets
V and W containing C \ U and D \ U respectively. Then (U V ) (U W ) = U,
so
(C, D) (UU0 V UU0 W ) I 0 (X) 1 (UU0 ).
If X is also T1 , it is a Hausdorff space, so compact sets are closed. Therefore
: K(X)
K(X)
K(X)
is continuous because its domain and range have the
0 (X) K
0 (X) and K
0 (X).
subspace topologies inherited from K
Let I(X) (resp. I 0 (X)) be the set of pairs (K, L) of compact (resp. closed)
subsets of X such that K L 6= , endowed with the topology it inherits from the
product topology of K(X) K(X) (resp. K0 (X) K0 (X)). The relevant topologies
are relative topologies obtained from the spaces in the last result, so:
Lemma 4.5.6. If X is a normal space, then : (C, D) 7 C D is a continuous
function from I 0 (X) to K0 (X). If, in addition, X is a T1 space, then : I(X)
K(X) is continuous.
4.5.3
Singletons
4.5.4
73
K}.
By Lemma 4.5.1, this establishes the asserted continuity when T {K,
1
To demonstrate continuity when T = H we must also show that (VW ) is open
in H(X) H(Y ) whenever W X Y is open. Suppose that (K L) W 6= .
Choose (x, y) (K L) W , and choose open neighborhoods U and V of x and y
with U V W . Then
K L VU VV 1 (VW ).
4.5.5
1
f (WV ) = Wf 1 (V ) for all W {U, U, V}.
There is the following consequence for closed sets.
K, H}, then f :
Lemma 4.5.11. If X is compact, Y is Hausdorff, and T {K,
K 7 f (K) is a continuous function from T 0 (X) to T 0 (Y ).
74
4.5.6
Whenever we have a set of subsets ofSsome space, we can take the union of its
elements. For any open U X we have KUU K = U because for each x U, {x}
is compact. Since the sets UU are a base for the topology of K(X), it follows that
the union of all elements of an open subset of K(X) is open. If U and V1 , . . . , Vk
are open, then UU VV1 VVk = if there is some j with U Vj = , and
otherwise
{x, y1 , . . . , yk } UU VV1 VVk
whenever x U and y1 V1 U, . . . , yk Vk U, so the union of all K
UU VV1 VVk is again U. Therefore the union of all the elements of an open
1
Proof: an open cover of the subset, together with its complement, is an open cover of the
space, any finite subcover of which yields a finite subcover of the subset.
2
Proof: fixing a point y in the complement of the compact set K, for each x K there are
disjoint neighborhoods of Ux of x and Vx of y, {Ux } is an open cover of K, and if Ux1 , . . . , Uxn is
a finite subcover, then Vx1 . . . Vxn is a neighborhood of y that does not intersect K.
75
subset of H(X) is open. If X is either T1 or regular, then similar logic shows that
for either T {K0 , H0 } the union of the elements of an open subset of T (X) is
open.
If a subset C of H(X) or H0 (X) is compact, then it is automatically compact in
the coarser topology of K(X) or K0 (X). Therefore the following two results imply
the analogous claims for the H(X) and H0 (X), which are already interesting.
S
Lemma 4.5.14. If S K(X) is compact, then L := KS K is compact.
Chapter 5
Topologies on Functions and
Correspondences
In order to study of robustness of fixed points, or sets of fixed points, with respect
to perturbations of the function or correspondence, one must specify topologies on
the relevant spaces of functions and correspondences. We do this by identifying
a function or correspondence with its graph, so that the topologies from the last
chapter can be invoked. The definitions of upper and lower semicontinuity, and their
basic properties, are given in Section 5.1. There are two topologies on the space of
upper semicontinuous correspondences from X to Y . The strong upper topology,
which is defined and discussed in Section 5.2, turns out to be rather poorly behaved,
and the weak upper topology, which is usually at least as coarse, is presented in
Section 5.3. When X is compact the strong upper topology coincides with the weak
upper topology.
We will frequently appeal to a perspective in which a homotopy h : X [0, 1]
Y is understood as a continuous function t 7 ht from [0, 1] to the space of continuous functions from X to Y . Section 5.4 presents the underlying principle in full
generality for correspondences. The specializations to functions of the strong and
weak upper topologies are known as the strong topology and the weak topology
respectively. If X is regular, then the weak topology coincides with the compactopen topology, and when X is compact the strong and weak topologies coincide.
Section 5.5 discusses these matters, and presents some results for functions that are
not consequences of more general results pertaining to correspondences.
The strong upper topology plays an important role in the development of the
topic, and its definition provides an important characterization of the weak upper
topology when the domain is compact, but it does not have any independent significance. Throughout the rest of the book, barring an explicit counterindication, the
space of upper semicontinuous correspondences from X to Y will be endowed with
the weak upper topology, and the space of continuous functions from X to Y will be
endowed with the weak topology.
76
5.1
77
5.2
Let X and Y be topological spaces with Y Hausdorff, and let U(X, Y ) be the set
of upper semicontinuous correspondences from X to Y . Proposition 5.1.2 insures
that the graph of each F U(X, Y ) is closed, so there is an embedding F 7 Gr(F )
of U(X, Y ) in K0 (X Y ). The strong upper topology is the topology induced
by this embedding when the image has the subspace topology. Let US (X, Y ) be
U(X, Y ) endowed with this topology. Since {UV0 : V X Y is open } is a subbase
for K0 (X Y ), there is a subbase of US (X, Y ) consisting of the sets of the form
{ F : Gr(F ) V }.
Naturally the following result is quite important.
Theorem 5.2.1. If Y is a Hausdorff space and X is a compact subset of Y , then
F P : US (X, Y ) K(X)
is continuous.
Proof. Since Y is Hausdorff, X and = { (x, x) : x X } are closed subsets of
Y and X Y respectively. For each F US (X, Y ), F P(F ) is the projection of
Gr(F ) onto the first coordinate. Since Gr(F ) is compact (Proposition 5.1.4) so
is Gr(F ), and the projection is continuous, so F P(F ) is compact. The definition
of the strong topology implies that Gr(F ) is a continuous function of F . Since is
closed in X Y , Lemma 4.5.3 implies that Gr(F ) is a continuous function of F ,
after which Lemma 4.5.10 implies that F P(F ) is a continuous function of F .
79
The basic operations for combining given correspondences to create new correspondences are restriction to a subset of the domain, cartesian products, and
composition. We now study the continuity of these constructions.
Lemma 5.2.2. If A is a closed subset of X, then the map F 7 F |A is continuous
as a function from US (X, Y ) to US (A, Y ).
Proof. Since A Y is a closed subset of X Y , continuity as a function from
US (X, Y ) to US (A, Y )that is, continuity of Gr(F ) 7 Gr(F ) (A Y )follows
immediately from Lemma 4.5.4.
An additional hypothesis is required to obtain continuity of restriction to a
compact subset of the domain, but in this case we obtain a kind of joint continuity.
Lemma 5.2.3. If X is regular, then the map (F, K) 7 Gr(F |K ) is a continuous
function from US (X, Y ) K(X) to K(X Y ). In particular, for any fixed K the
map F 7 F |K is a continuous function from US (X, Y ) to US (K, Y ).
Proof. Fix F US (X, Y ), K K(X), and an open neighborhood W of Gr(F |K ).
For each x K Lemma 4.5.8 gives neighborhoods Ux of x and Vx of F (x) with
Ux Vx W . Choose x1 , . . . , xk such that U := Ux1 . . . Uxk contains K.
Since X is regular, each point in K has a closed neighborhood contained in U, and
the interiors of finitely many of these cover K, so K has a closed neighborhood C
contained in U. Let
W := (Ux1 Vx1 ) . . . (Uxk Vxk ) ((X \ C) Y ).
Then (K, Gr(F )) Uint C UW , and whenever (K , Gr(F )) Uint C UW we have
Gr(F |K ) W (C Y ) (Ux1 Vx1 ) . . . (Uxk Vxk ) W.
Let X and Y be two other topological spaces with Y Hausdorff. Since the map
(C, D) 7 C D is not a continuous operation on closed sets, we should not expect
the function (F, F ) 7 F F from US (X, Y )US (X , Y ) to US (XX , Y Y ) to be
continuous, and indeed, after giving the matter a bit of thought, the reader should
be able to construct a neighborhood of the graph of the function (x, x ) 7 (0, 0)
that shows that the map (F, F ) 7 F F from US (R, R) US (R, R) to US (R2 , R2 )
is not continuous.
We now turn our attention to composition. Suppose that, in addition to X and
Y , we have a third topological space Z that is Hausdorff. (We continue to assume
that Y is Hausdorff.) We can define a composition operation from (F, G) 7 G F
from U(X, Y ) U(Y, Z) to U(X, Z) by letting
[
G(F (x)) :=
F (y).
yF (x)
That is, G(F (x)) is the projection onto Z of Gr(G|F (x)), which is compact by
Proposition 5.1.4, so G(F (x)) is compact. Thus G F is compact valued. To show
5.3
As in the last section, X and Y are topological spaces with Y Hausdorff. There
is another topology on U(X, Y ) that is in certain ways more natural and better
behaved than the strong upper topology. Recall that if {Bi }iI is a collection of
topological spaces and { fi : A Bi }iI is a collection of functions, the quotient
topology on A induced by this data is the coarsest topology such that each fi
is continuous. The weak upper topology on U(X, Y ) is the quotient topology
induced by the functions F 7 F |K US (K, Y ) for compact K X. Since a
function is continuous if and only if the preimage of every subbasic subset of the
range is open, a subbase for the weak upper topology is given by the sets of the
form { F : Gr(F |K ) V } where K X is compact and V is a (relatively) open
subset of K Y .
Let UW (X, Y ) be U(X, Y ) endowed with the weak upper topology. As in the
last section, we study the continuity of basic constructions.
Lemma 5.3.1. If A is a closed subset of X, then the map F 7 F |A is continuous
as a function from UW (X, Y ) to UW (A, Y ).
81
Proof. We need to show that the identity map from US (X, Y ) to UW (X, Y ) is continuous, which is to say that for any given compact K X, the map Gr(F )
Gr(F |K ) = Gr(F ) (K Y ) is continuous. This follows from Lemma 5.3.1 because
K Y is closed in X Y whenever K is compact.
If X is compact, the continuity of the identity map from UW (X, Y ) to US (X, Y )
follows directly from the definition of the weak upper topology.
There is a useful variant of Lemma 5.2.3.
Lemma 5.3.3. If X is normal, Hausdorff, and locally compact, then the function
(K, F ) 7 Gr(F |K ) is a continuous function from K(X) UW (X, Y ) to K(X Y ).
Proof. We will demonstrate continuity at a given point (K, F ) in the domain. Local
compactness implies that there is a compact neighborhood C of K. The map
F 7 F |C from U(X, Y ) to US (C, Y ) is a continuous function by virtue of the
definition of the topology of U(X, Y ). Therefore Lemma 5.2.3 implies that the
composition (K , F ) (K , F |C ) Gr(F |K ) is continuous, and of course it
agrees with the function in question on a neighborhood of (K, F ).
In contrast with the strong upper topology, for the weak upper topology cartesian products and composition are well behaved. Let X and Y be two other spaces
with Y Hausdorff.
Lemma 5.3.4. If X and X are Hausdorff, then the function (F, F ) 7 F F
from UW (X, Y ) UW (X , Y ) to UW (X X , Y Y ) is continuous.
Proof. First suppose that X and X are compact. Then, by Proposition 5.1.4,
the graphs of upper semicontinuous functions with these domains are compact,
and continuity of the function (F, F ) 7 F F from US (X, Y ) US (X , Y ) to
US (X X , Y Y ) follows from Proposition 4.5.9.
Because UW (X X , Y Y ) has the quotient topology, to establish the general case we need to show that (F, F ) 7 F F |C is a continuous function from
UW (X, Y ) UW (X , Y ) to US (C, Y Y ) whenever C X X is compact. Let
K and K be the projections of C onto X and X respectively; of course these sets
are compact. The map in question is the composition
(F, F ) (F |K , F |K ) F |K F |K (F |K F |K )|C .
The continuity of the second map has already been established, and the continuity
of the first and third maps follows from Lemma 5.3.1, because compact subsets of
Hausdorff spaces are closed and products of Hausdorff spaces are Hausdorff1 .
Suppose that, in addition to X and Y , we have a third topological space Z that
is Hausdorff.
Lemma 5.3.5. If K X is compact, Y is normal and locally compact, and X
Y Z is normal, then
(F, G) 7 Gr(G F |K )
is a continuous function from UW (X, Y ) UW (Y, Z) to K(X Z).
1
I do not know if the compact subsets of X X are closed when X and X are compact spaces
whose compact subsets are closed.
5.4
83
5.5
Continuous Functions
Chapter 6
Metric Space Theory
In this chapter we develop some advanced results concerning metric spaces.
An important tool, partitions of unity, exist for locally finite open covers of
a normal space: this is shown in Section 6.2. But sometimes we will be given
a local cover that is not necessarily locally finite, so we need to know that any
open cover has a locally finite refinement. A space is paracompact if this is the
case. Paracompactess is studied in Section 6.1; the fact that metric spaces are
paracompact will be quite important.
Section 6.3 describes most of the rather small amount we will need to know
about topological vector spaces. Of these, the most important for us are the locally
convex spaces, which have many desirable properties. One of the larger themes of
this study is that the concepts and results of fixed point theory extend naturally to
this level of generality, but not further.
Two important types of topological vector spaces, Banach spaces and Hilbert
spaces, are introduced in Section 6.4. Results showing that metric spaces can be
embedded in such linear spaces are given in Section 6.5. Section 6.6 presents an
infinite dimensional generalization of the Tietze extension theorem due to Dugundji.
6.1
Paracompactness
86
For each x there is a least such that x U and an n large enough that (c) holds,
so x Vn unless x Vj for some and j < n. Thus {Vn } is a cover of X, and
of course each Vn is open and contained in U , so it is a refinement of {U }.
To prove that the cover is locally finite we fix x, let be the least element of
A such that x Vn for some n, and choose j such that U2j (x) Vn . We claim
that U2nj (x) intersects only finitely many Vi .
If i > j and y satisfies (a)-(c) with and i in place of and n, then U2nj (x)
/ Vn , and n + j, i j + 1. Therefore
U2i (y) = because U2j (x) Vn , y
U2nj (x) Vi = .
For i j we will show that there is at most one such that U2nj (x) intersects
Vi . Suppose that y and z are points satisfying (a)-(c) for and , with i in place
of j. Without loss of generality preceeds . Then U32i (y) U , z
/ U , and
n + j > i, so U2nj (x) cannot intersect both U2i (y) and U2i (z). Since this is
the case for all y and z, U2nj (x) cannot intersect both Vi and Vi .
6.2
Partitions of Unity
87
Theorem 6.2.2. For any locally finite open cover {U }A of a normal space X
there is a partition of unity subordinate to {U }.
A basic tool used in the constructive proof of this result, and many others, is:
Lemma 6.2.3 (Urysohns Lemma). If X is a normal space and C U X with C
closed and U open, then there is a continuous function : X [0, 1] with (x) = 0
for all x C and (x) = 1 for all x X \ U.
Proof. Since X is normal, whenever C U , with C closed and U open, there
exist a closed C and an open U such that C U , X \ U X \ C , and
U (X \ C ) = , which is to say that C U C U . Let C0 := C and
U1 := U. Choose an open U1/2 and a closed C1/2 with C0 U1/2 C1/2 U1 .
Choose an open U1/4 and a closed C1/4 with C0 U1/4 C1/4 U1/2 , and choose
an open U3/4 and a closed C3/4 with C1/2 U3/4 C3/4 U1 . Continuing in
this fashion, we obtain a system of open sets Ur and a system of closed sets Cr for
rationals r [0, 1] of the form k/2m (except that C1 and U0 are undefined) with
Ur Cr Us Cs whenever r < s.
For x X let
(
S
inf{ r : x Cr }, x r Cr
(x) :=
1,
otherwise.
Clearly (x) = 0 for all x C and (x) = 1 for all x X \ U. Any open subset
of [0, 1] is a union of finite intersections of sets of the form [0, a) and (b, 1], where
0 < a, b < 1, and
[
[
1 [0, a) =
Ur and 1 (b, 1] = (X \ Cr )
r<a
r>b
Our goal is to find such an F with B = A. The partial thinnings can be partially
ordered as follows: F < G if the domain of F is a proper subset of the domain of
G and F and G agree on this set. We will show that this ordering has maximal
elements, and that the domain of a maximal element is all of A.
88
A\B,6=
:= P
6.3
89
90
6.4
91
0 hv, viw hv, wiv, hv, viw hv, wiv = hv, vi hv, vihw, wi hv, wi2 ,
which implies the Cauchy-Schwartz inequality: hv, wi kvk kwk for all v, w
V . This holds with equality if v = 0 or hv, viw hv, wiv, which is the case if
and only if w is a scalar multiple of v, and otherwise the inequality is strict. The
Cauchy-Schwartz inequality implies the inequality in the calculation
kv + wk2 = hv + w, v + wi = kvk2 + 2hv, wi + kwk2 (kvk + kwk)2 ,
which implies (c) and completes the verification and k k is a norm. A vector space
endowed with an inner product and the associated norm and topology is called an
inner product space. A Hilbert space is a complete inner product space.
Up to linear isometry there is only one separable2 Hilbert space. Let
H := { s = (s1 , s2 , . . .) R : s21 + s22 + < }
2
Recall that a metric space is separable if it contains a countable set of points whose closure
is the entire space.
92
P
be the Hilbert space of square summable sequences. Let hs, ti := i si ti be the usual
inner product; the Cauchy-Schwartz inequality implies that this sum is convergent.
For any Cauchy sequence in H and for each i, the sequence of ith components is
Cauchy, and the element of R whose ith component is the limit of this sequence is
easily shown to be the limit in H of the given sequence. Thus H is complete. The
set of points with only finitely many nonzero components, all of which are rational,
is a countable dense subset, so H is separable.
We wish to show that any separable Hilbert space is linearly isomorphic to H, so
let V be a separable Hilbert space, and let {v1 , v2 , . . . } be a countable dense subset.
The span of this set is also dense, of course. Using the Gram-Schmidt process, we
may pass from this set to a countable sequence w1 , w2, . . . of orthnormal vectors
that has the same span. It is now easy to show that s 7 s1 w1 + s2 w2 + is a
linear isometry between H and V .
6.5
EmbeddingTheorems
93
Pk
i=1 i fyi
6.6
Dugundjis Theorem
The well known Tietze extension theorem asserts that if a topological space X
is normal and f : A [0, 1] is continuous, where A X is closed, then f has a
continuous extension to all of X. A map into a finite dimensional Euclidean space
is continuous if its component functions are each continuous, so Tietzes theorem
is adequate for finite dimensional applications. Mostly, however, we will work with
spaces that are potentially infinite dimensional, for which we will need the following
variant due to Dugundji (1951).
Theorem 6.6.1. If A is a closed subset of a metric space (X, d), Y is a locally
convex topological vector space, and f : A Y is continuous, then there is a
continuous extension f : X Y whose image is contained in the convex hull of
f (A).
94
Proof. The sets Ud(x,A)/2 (x) are open and cover X \ A. Theorem 6.1.1 implies the
existence of an open locally finite refinement {W }I . Theorem 6.2.2 implies the
existence of a partition of unity { }I subordinate to {W }I . For each choose
a A with d(a , W ) < 2d(A, W ), and define the extension by setting
X
f (x) :=
(x)f (a ) (x X \ A).
I
Chapter 7
Retracts
This chapter begins with Kinoshitas example of a compact contractible space
that does not have the fixed point property. The example is elegant, but also rather
complex, and nothing later depends on it, so it can be postponed until the reader is
in the mood for a mathematical treat. The point is that fixed point theory depends
on some additional condition over and above compactness and contractibility.
After that we develop the required material from the theory of retracts. We
first describe retracts in general, and then briefly discuss Euclidean neighborhood
retracts, which are retracts of open subsets of Euclidean spaces. This concept is
quite general, encompassing simplicial complexes and (as we will see later) smooth
manifolds.
The central concept of the chapter is the notion of an absolute neighborhood
retract (ANR) which is a metrizable space whose image, under any embedding as a
closed subset of a metric space, is a retract of some neighborhood of itself. The two
key characterization results are that an open subset of a convex subset of a locally
convex linear space is an absolute neighborhood retract, and that an ANR can be
embedded in a normed linear space as a retract of an open subset of a convex set.
An absolute retract (AR) is a space that is a retract of any metric space it
is embedded in as a closed subset. It turns out that the ARs are precisely the
contractible ANRs.
The extension of fixed point theory to infinite dimensional settings ultimately
depends on approximating the setting with finite dimensional objects. Section
7.6 provides one of the key results in this direction.
7.1
Kinoshitas Example
This example came to be known as the tin can with a roll of toilet paper. As
you will see, this description is apt, but does not do justice to the examples beauty
and ingenuity.
Polar coordinates facilitate the description. Let P = [0, ) R, with (r, ) P
identified with the point (r cos , r sin ). The unit circle and the open unit disk are
C = { (r, ) : r = 1 }
and
95
D = { (r, t) : r < 1 }.
96
CHAPTER 7. RETRACTS
Then S is a curve that spirals out from the origin, approaching C asymptotically.
The space of the example is
X = (C [0, 1]) (D {0}) (S [0, 1]) R3 .
Here C [0, 1] is the cylindrical side of the tin can, D {0} is its base, and S [0, 1]
is the roll of toilet paper. Evidently X is closed, hence compact, and there is an
obvious contraction of X that first pushes the cylinder of the tin can and the toilet
paper down onto the closed unit disk and then contracts the disk to the origin.
We are now going to define functions
f1 : C [0, 1] X,
f2 : D {0} X,
f3 : S [0, 1] X
7.2. RETRACTS
97
7.2
Retracts
This section prepares for later material by presenting general facts about retractions and retracts. Let X be a metric space, and let A be a subset of X such that
there is a continuous function r : X A with r(a) = a for all a A. We say that
A is a retract of X and that r is a retraction. Many desirable properties that X
might have are inherited by A.
Lemma 7.2.1. If X has the fixed point property, then A has the fixed point property.
Proof. If f : A A is continuous, then f r necessarily has a fixed point, say a ,
which must be in A, so that a = f (r(a )) = f (a ) is also a fixed point of f .
Lemma 7.2.2. If X is contractible, then A is contractible.
Proof. If c : X [0, 1] X is a contraction X, then so is (a, t) 7 r(c(a, t)).
Lemma 7.2.3. If X is connected, then A is connected.
Proof. We show that if A is not connected, then X is not connected. If U1 and
U2 are nonempty open subsets of A with U1 U2 = and U1 U2 = A, then
r 1 (U1 ) and r 1 (U2 ) are nonempty open subsets of X with r 1 (U1 ) r 1 (U2 ) =
and r 1 (U1 ) r 1 (U2 ) = X.
Here are two basic observations that are too obvious to prove.
Lemma 7.2.4. If s : A B is a second retraction, then s r : X B is a
retraction, so B is a retract of X.
Lemma 7.2.5. If A Y X, then the restriction of r to Y is a retraction, so A
is a retract of Y .
98
CHAPTER 7. RETRACTS
7.3
99
100
CHAPTER 7. RETRACTS
Proof. Suppose that such an r exists, and let g : D n S n1 be the function that
takes each x S n1 to itself and takes each x D n \ S n1 to the point where the
line segment between r(x) and x intersects S n1 . An easy argument shows that g
is continuous at each x D n \ S n1 , and another easy argument shows that g is
continuous at each x S n1 , so g is continuous. If a : S n1 S n1 is the antipodal
map a(x) = x, then a g gives a map from D n to itself that does not have a fixed
point, contradicting Brouwers fixed point theorem.
7.4
101
102
CHAPTER 7. RETRACTS
7.5
Absolute Retracts
103
104
CHAPTER 7. RETRACTS
Proof of Proposition 7.5.4. Let Y , f : X Y , and h : A [0, 1] Y satisfy (a)(d) above. By Theorem 6.5.2 we may assume without loss of generality that Y is
contained in a Banach space S, and is a relatively closed subset of its convex hull C.
Let Z := (X {0})(A[0, 1]), and define g : Z Y by setting g(x, 0) = f (x) and
g(a, t) = h(a, t). Dugundjis theorem implies that there is a continuous extension
g : X [0, 1] C of g. Let W C be a neighborhood of Y for which there is a
retraction r : W Y , let V := g 1 (W ), and let : X [0, 1] V be a continuous
map that is the identity on Z, as per the result above. Clearly := r g has
the indicated properties.
We now return to the characterization of ARs.
Proof of Proposition 7.5.3. Let A be an ANR. By Theorem 6.5.2 we may embed A
as a relatively closed subset of a convex subset C of a Banach space.
If A is an AR, then it is a retract of C. A convex set is contractible, and a
retract of a contractible set is contractible (Lemma 7.2.2) so A is contractible.
Suppose that A is conractible. By Proposition 7.5.1 it suffices to show that A is
a retract of C. Let c : A[0, 1] A be a contraction, and let a1 be the final value
a1 , by which we mean that c(a, 1) = a1 for all a A. Set Z := (C {0})(A[0, 1]),
and define f : Z A by setting f (x, 0) := a1 for x C and f (a, t) := c(a, 1 t) for
(a, t) A [0, 1]. Proposition 7.5.4 implies the existence of a continuous extension
f : C [0, 1] A. Now r := f (, 1) : C A is the desired retraction.
7.6
Domination
In our development of the fixed point index an important idea will be to pass
from a theory for certain simple or elementary spaces to a theory for more general
spaces by showing that every space of the latter type can be approximated by a
simpler space, in the sense of the following definitions. Fix a metric space (X, d).
Definition 7.6.1. If Y is a topological space and > 0, a homotopy : Y [0, 1]
X is an -homotopy if
d (y, s), (y, t) <
for all y Y and all 0 s, t 1. We say that 0 and 1 are -homotopic.
105
7.6. DOMINATION
Vj1 ...Vjk 6=
Of course the denominator is always positive, so these functions are well defined
and continuous. There is a continuous function : C N(U1 ,...,Un ) given by
(x) :=
n
X
j (x)ej .
j=1
n
X
j=1
j ej = r
n
X
j=1
j xj .
Pn
Consider a point y =
j=1 j ej N(U1 ,...,Un ) . Let j1 , . . . , jk be the indices j
such that j > 0, ordered so that (xj1 ) max{(xj2 ), . . . , (xjk )}. Let
Tp B :=
U2(xj1 ) (xj1 ). The definition of N(U1 ,...,Un ) implies that there is a point z h=0 Ujh .
For all h = 1, . . . , k we have xjh B because
d(z, xjh ) < (xjh ) (xj1 ).
Now note that
B r 1 (U/2 (xj1 ) X) U.
P
Since B is convex, it contains kh=1 jh xjh , so is well defined.
Now we would like to define a homotopy : C [0, 1] X by setting
X
(x, t) = r (1 t)
j (x)xj + tx ,
j
106
CHAPTER 7. RETRACTS
Chapter 8
Essential Sets of Fixed Points
Figure 2.1 shows a function f : [0, 1] [0, 1] with two fixed points, s and t.
Intuitively, they are qualitatively different, in that a small perturbation of f can
result in a function that has no fixed points near s, but this is not the case for t.
This distinction was recognized by Fort (1950) who described s as inessential, while
t is said to be essential.
1
b
Figure 1.1
In game theory one often deals with correspondences with sets of fixed points
that are infinite, and include continua such as submanifolds. As we will see, the
definition proposed by Fort can be extended to sets of fixed points rather easily:
roughly, a set of fixed points is essential if every neighborhood of it contains fixed
points of every sufficiently close perturbation of the given correspondence. (Here
one needs to be careful, because in the standard terminology of game theory, following Jiang (1963), essential Nash equilibria, and essential sets of Nash equilibria,
are defined in terms of perturbations of the payoffs. This is a form of Q-robustness,
which is studied in Section 8.3.) But it is easy to show that the set of all fixed
107
108
8.1
We now extend the Kakutani fixed point theorem to correspondences with infinite dimensional domains. The result below was proved independently by Fan
(1952) and Glicksberg (1952) using quite similar methods; our proof is perhaps a
bit closer to Fans. In a sense the result was already known, since it can be derived
from the Eilenberg-Montgomery theorem, but the proof below is much simpler.
Theorem 8.1.1 (Fan, Glicksberg). If V is a locally convex topological vector space,
X V is nonempty, convex, and compact, and F : X X is an upper semicontinuous convex valued correspondence, then F has a fixed point.
We treat two technical points separately:
Lemma 8.1.2. If V is a (not necessarily locally convex) topological vector space
and K, C V with K compact and C closed, then K + C is closed.
109
Proof. We will show that the compliment is open. Let y be a point of V that is not
in K + C. For each x K, translation invariance of the topology of V implies that
x + C is closed, so Lemma 6.3.2 gives a neighborhood Wx of the origin such that
(y + Wx + Wx ) (x + C) = . Since we can replace Wx with Wx Wx , we may
assume that Wx = Wx , so that (y + Wx ) (x + C + Wx ) = . Choose x1 , . . . , xk
such that the sets xi + Wxi cover K, and let W = Wx1 . . . Wxk . Now
[
(y + W ) (K + C) (y + W ) (xi + C + Wxi )
i
[
i
110
8.2
111
()
()
112
8.3
Kinoshitas Theorem
113
114
Proof. The continuity of F P (Theorem 5.2.1) implies that for any neighborhood
V X of F P(F ) there is a neighborhood U A of a0 such that F P(Q(a)) V
for all a U. The Fan-Glicksberg fixed point theorem implies that F P(Q(a)) is
nonempty.
This result shows that if our goal is to discriminate between some fixed points
and others, these concepts must be strengthened in some way. The two main
methods for doing this are to require either connectedness or minimality.
Definition 8.3.7. A nonempty compact set K F P(F ) is a minimal Q-robust
set if it is Q-robust and minimal in the class of such sets: K is Q-robust and no
proper subset is Q-robust. A minimal connected Q-robust set is a connected
Q-robust set that does not contain a proper subset that is connected and Q-robust.
In general a minimal Q-robust set need not be connected. For example, if
(A, a0 ) = ((1, 1), 0) and Q(a)(t) = argmaxt[0,1] at (so that F (t) = [0, 1] for all t)
then F P(Q(a)) is {0} if a < 0 and it is {1} if a > 0, so the only minimal Q-robust
set is {0, 1}. In view of this one must be careful to distinguish between a minimal
connected Q-robust set and a minimal Q-robust set that happens to be connected.
Theorem 8.3.8. If K F P(F ) is a Q-robust set, then it contains a minimal
Q-robust set, and if K is a connected Q-robust set, then it contains a minimal
connected Q-robust set.
Proof. Let C be the set of Q-robust sets that are contained in K. We order this set
by reverse inclusion, so that our goal is to show that C has a maximal element. This
follows from Zorns lemma if we can show that any completely ordered subset O has
an upper bound in C. The finite intersection property implies that the intersection
of all elements of O is nonempty; let K be this intersection. If K is not Qrobust, then there is a neighborhood V of K such that every neighborhood U of
a0 contains a point a such that Q(a) has no fixed points in V . If L O, we cannot
have L V because L is Q-robust, but now { L \ V : L O } is a collection of
compact sets with the finite intersection property, so it has a nonempty intersection
that is contained in K but disjoint from V . Of course this is absurd.
The argument for connected Q-robust sets follows the same lines, except that in
addition to showing that K is Q-robust, we must also show that it is connected.
If not there are disjoint open sets V1 and V2 such that K V1 V2 and K V1 6=
=
6 K V2 . For each L O we have L V1 6= =
6 L V2 , so L \ (V1 V2 ) must
be nonempty because L is connected. As above, { L \ (V1 V2 ) : L O } has a
nonempty intersection that is contained in K but disjoint from V1 V2 , which is
impossible.
Chapter 9
Approximation of
Correspondences
In extending fixed point theory from functions to correspondences, an important
method is to show that continuous functions are dense in the space of correspondences, so that any correspondence can be approximated by a function. In the
last chapter we saw such a result (Theorem 8.2.1) for convex valued correspondences, but much greater care and ingenuity is required by the arguments showing
that contractible valued correspondences have good approximations. This chapter states and proves the key result in this direction. This result was proved in
the Euclidean case by Mas-Colell (1974) and extended to ANRs by the author in
McLennan (1991).
9.1
Our main result can be stated rather easily. We now fix ANRs X and Y . We
assume throughout this chapter that X is separable, in order to be able to invoke
the domination theorem.
Theorem 9.1.1 (Approximation Theorem). Suppose that C and D are compact
subsets of X with C int D. Let F : D Y be an upper semicontinuous contractible valued correspondence. Then for any neighborhood U of Gr(F |C ) there
are:
(a) a continuous f : C Z with Gr(f ) U;
(b) a neighborhood U of Gr(F ) such that, for any two continuous functions f0 , f1 :
D Y with Gr(f0 ), Gr(f1 ) U , there is a homotopy h : C [0, 1] Y with
h0 = f0 |C , h1 = f1 |C , and Gr(ht ) U for all 0 t 1.
Roughly, (a) is an existence result, while (b) is uniqueness up to effective equivalence.
Here, and later in the book, things would be much simpler if we could have
C = D. More precisely, it would be nice to drop the assumption that C int D.
This may be possible (that is, I do not know a relevant counterexample) but a proof
would certainly involve quite different methods.
115
116
9.2
The proof of Theorem 9.1.1 begins with a concrete geometric construction that
is given in this section. In subsequent sections we will transport this result to
increasingly general settings, eventually arriving at our objective.
We now fix a locally convex topological vector space T and a convex Q T . A
subset Z of a vector space is balanced if z Z whenever z Z and || 1. Since
T is locally convex, every neighborhood of the origin contains a convex neighborhood
U, and U U is a neighborhood that is convex and balanced. Working with
balanced neighborhoods of the origin allows us to not keep track of the difference
between a neighborhood and its negation.
Proposition 9.2.1. Let A and B be convex balanced neighborhoods of the origin in
T with 2A B. Suppose S Q is compact and c : S[0, 1] S is a contraction for
which there is a > 0 such that c(s, t) c(s , t ) B for all (s, t), (s, t ) S [0, 1]
with s s 3A and |t t | < . Let L be a simplex. Then any continuous
f : L (S + A) Q has a continuous extension f : L (S + B) Q.
Proof. Let be the barycenter of L. We define polar coordinate functions
y : L \ {} L and t : L \ {} [0, 1)
implicitly by requiring that
(1 t(x))y(x) + t(x) = x.
Let
L1 = t1 ([0, 31 ]),
L2 = t1 ([ 13 , 32 ]),
L3 = t1 ([ 32 , 1)) {}.
117
Let d be a metric on L. Since f , t(), and y() are continuous, and L2 is compact,
for some sufficiently small > 0 it is the case that
f (y(x)) f (y(x)) A and |t(x) t(x )| < 31
for all x, x L2 such that d(x, x ) < . There is a polyhedral subdivision of L2
whose cells are the sets
y 1 (F ) t1 ( 31 ),
y 1 (F ) L2 ,
y 1(F ) t1 ( 32 )
for the various faces F of L. Proposition 2.5.2 implies that repeated barycentric
subdivision of this polyhedral complex results eventually in a simplicial subdivision
of L2 whose mesh is less than .
For each vertex v of this subdivision choose s(v) (f (y(v)) + A) S, and set
f (v) := c(s(v), 3t(v) 1).
Let be a simplex of the subdivision of L2 with vertices v1 , . . . , vr . We define f
on by linear interpolation on : if x = 1 v1 + + r vr , then
f (x) := 1 f (v1 ) + + r f (vr ).
This definition does not depend on the choice of if x is contained in more than
one simplex, it is continuous on each , and the simplices are a finite closed cover
of L2 , so f is continuous.
Suppose that v and v are two vertices of , so they are the endpoints of an
edge. We have d(v, v ) < , so f (y(v)) f (y(v )) A and |t(v) t(v )| < 31 . In
addition, s(v) f (y(v)) and f (y(v )) s(v ) are elements of A, so
s(v) s(v ) 3A and |(3t(v) 1) (3t(v ) 1)| < ,
from which it follows, by hypothesis, that f (v) f (v ) B. Consider a point
x = 1 v1 + + r vr . Since f (v1 ) S and
f (x) f (v1 ) =
r
X
j=1
j (f (vj ) f (v1 ))
118
The point s(vj ) was chosen with f (y(vj )) s(vj ) A, and f (y(x)) f (y(vj )) A
because d( 32 y(x) + 13 , y(vj )) < , so
f (y(x)) (s(vj ) + 2A) Q (s(vj ) + B) Q.
Since f (x) is a convex combination of f (y(x)) and f ( 32 y(x) + 13 ) we have
f (x) (s(vj ) + B) Q (S + B) Q.
Thus f (L1 ) (S + B) Q.
Let z be the point S is contracted to by c: c(S, 1) = {z}. We define f on L3 by
setting f (x) := z. Of course this is a continuous function whose image is contained
in S (S + B) Q.
If x L1 L2 , then t(x) = 31 and 23 y(x) + 31 = x, so the formula defining f on
L1 agrees with the definition of f for elements of L2 at x. If v is a vertex of the
subdivision of L2 contained in L2 L3 , then t(v) = 23 , so that the definition of f on
L2 gives f (v) = c(s(v), 3t(v) 1) = z. If x L2 L3 , then L2 L3 contains any
simplex of the subdivision of L2 that has x as an element, and the definition of f
on L2 gives f (x) = z. Thus this definition agrees with the definition of f on L2 at
points in L2 L3 . Thus f is well defined and continuous.
9.3
119
and let Vy be an open subset of U(52)ry (y) that contains y. Then for all y, y X,
if Vy Vy 6= , then Vy Ury (y).
120
p
[
j=1
j J
9.4
The next step is a result in which the domains are subsets of the ANR X.
Proposition 9.4.1. Suppose that C D X where C and D are compact with
C int D. Let F : D Z be an upper semicontinuous contractible valued correspondence. Then for any neighborhood V of Gr(F |C ) there exist:
(a) a continuous f : C Z with Gr(f ) V ;
121
h1 = i f1 |C ,
so that
r h0 = f0 |C ,
r h1 = f1 |C ,
Proof. For each x C Lemma 9.3.2 allows us to choose x > 0 and a neighborhood
Ax of F (x) such that Ux (x) Ax V . Replacing x with a smaller number
if need be, we may assume without loss of generality that F (x ) Ax for all
x Ux (x). Choose x1 , . . . , xH such that Ux1 /2 (x1 ), . . . , UxH /2 (xH ) cover C. Let
:= min{xi /2}, and set
[
V :=
Uxi /2 (xi ) Axi .
i
122
123
Part II
Smooth Methods
124
Chapter 10
Differentiable Manifolds
This chapter introduces the basic concepts of differential topology: manifold,
tangent vector, smooth map, derivative. If these concepts are new to you,
you will probably be relieved to learn that these are just the basic concepts of
multivariate differential calculus, with a critical difference.
In multivariate calculus you are handed a coordinate system, and a geometry,
when you walk in the door, and everything is a calculation within that given Euclidean space. But many of the applications of multivariate calculus take place in
spaces like the sphere, or the physical universe, whose geometry is not Euclidean.
The theory of manifolds provides a language for the concepts of differential calculus
that is in many ways more natural, because it does not presume a Euclidean setting.
Roughly, this has two aspects:
In differential topology spaces that are locally homeomorphic to Euclidean
spaces are defined, and we then impose structure that allows us to talk about
differentiation of functions between such spaces. The concepts of interest
to differential topology per se are those that are invariant under diffeomorphism, much as topology is sometimes defined as rubber sheet geometry,
namely the study of those properties of spaces that dont change when the
space is bent or stretched.
The second step is to impose local notions of angle and distance at each point
of a manifold. With this additional structure the entire range of geometric
issues can be addressed. This vast subject is called differential geometry.
For us differential topology will be primarily a tool that we will use to set up an
environment in which issues related to fixed points have a particularly simple and
tractable structure. We will only scratch its surface, and differential geometry will
not figure in our work at all.
The aim of this chapter is provide only as much information as we will need
later, in the simplest and most concrete manner possible. Thus our treatment of
the subject is in various ways terse and incomplete, even as an introduction to this
topic, which has had an important influence on economic theory. Milnor (1965) and
Guillemin and Pollack (1974) are recommended to those who would like to learn a
bit more, and at a somewhat higher level Hirsch (1976) is more comprehensive, but
still quite accessible.
125
126
10.1
We begin with a quick review of the most important facts of multivariate differential calculus. Let f : U Rn be a function where U Rm is open. Recall
that if r 1 is an integer, we say that f is C r if all partial derivatives of order
r are defined and continuous. For reasons that will become evident in the next
paragraph, it can be useful to extend this notation to include r = 0, with C 0 interpreted as a synonym for continuous. We say that f is C if it is C r for all finite
r. An order of differentiability is either a nonnegative integer r or , and we
write 2 r , for example, to indicate that r is such an object, within the given
bounds.
If f is C 1 , then f is differentiable: for each x U and > 0 there is > 0
such that
kf (x ) f (x) Df (x)(x x)k kx xk
for all x U with kx xk < , where the derivative of f at x is the linear
function
Df (x) : Rm Rn
given by the matrix of first partial derivatives at x. If f is C r , then the function
Df : U L(Rm , Rn )
is C r1 if we identify L(Rm , Rn ) with the space Rnm of n m matrices. The
reader is expected to know the standard facts of elementary calculus, especially
that addition and multiplication are C , so that functions built up from these
operations (e.g., linear functions and matrix multiplication) are known to be C .
There are three basic operations used to construct new C r functions from give
functions. The first is restriction of the function to an open subset of its domain,
which requires no comment because the derivative is unaffected. The second is
forming the cartesian product of two functions: if f1 : U Rn1 and f2 : U Rn2
are functions, we define f1 f2 : U Rn1 +n2 to be the function x 7 (f1 (x), f2 (x)).
Evidently f1 f2 is C r if and only if f1 and f2 are C r , and when this is the case
we have
D(f1 f2 ) = Df1 Df2 .
The third operation is composition. The most important theorem of multivariate
calculus is the chain rule: if U Rm and V Rn are open and f : U V and
g : V Rp are C 1 , then g f is C 1 and
D(g f )(x) = Dg(f (x)) Df (x)
for all x U. Of course the composition of two C 0 functions is C 0 . Arguing inductively, suppose we have already shown that the composition of two C r1 functions
is C r1 . If f and g are C r , then Dg f is C r1 , and we can apply the result above
about cartesian products, then the chain rule, to the composition
x 7 (Dg(f (x)), Df (x)) 7 Dg(f (x)) Df (x)
to show that D(g f ) is C r1 , so that g f is C r .
127
Often the domain and range of the pertinent functions are presented to us as
vector spaces without a given or preferred coordinate system, so it is important to
observe that we can use the chain rule to achieve definitions that are independent of
the coordinate systems. Let X and Y be m- and n-dimensional vector spaces. (In
this chapter all vector spaces are finite dimensional, with R as the field of scalars.)
Let c : X Rm and d : Y Rn be linear isomorphisms. If U X is open, we can
say that a function f : U Y is C r , by definition, if d f c1 : c(U) Rk is C r ,
and if this is the case and x U, then we can define the derivative of f at x to be
Df (x) = d1 D(d f c1 )(c(x)) c L(X, Y ).
Using the chain rule, one can easily verify that these definitions do not depend
on the choice of c and d. In addition, the chain rule given above can be used to
show that this coordinate free definition also satisfies a chain rule. Let Z be a
third p-dimensional vector space. Then if V Y is open, g : V Z is C r , and
f (U) V , then g f is C r and D(g f ) = Dg Df .
Sometimes we will deal with functions whose domains are not open, and we need
to define what it means for such a function to be C r . Let S be a subset of X of any
sort whatsoever. If Y is another vector space and f : S Y is a function, then
f is C r by definition if there is an open U X containing S and a C r function
F : U Y such that f = F |S . Evidently being C r isnt the same thing as having
a well defined derivative at each point in the domain!
Note that the identity function on S is always C r , and the chain rule implies
that compositions of C r functions are C r . Those who are familiar with the category
concept will recognize that there is a category of subsets of finite dimensional vector
spaces and C r maps between them. (If you havent heard of categories it would
certainly be a good idea to learn a bit about them, but what happens later wont
depend on this language.)
We now state coordinate free versions of the inverse and implicit function theorems. Since you are expected to know the usual, coordinate dependent, formulations
of these results, and it is obvious that these imply the statements below, we give
no proofs.
Theorem 10.1.1 (Inverse Function Theorem). If n = m (that is, X and Y are
both m-dimensional) U X is open, f : U Y is C r , x U, and Df (x) is
nonsingular, then there is an open V U containing x such that f |V is injective,
f (V ) is open in Y , and (f |V )1 is C r .
Suppose that U X Y is open and f : U Z is a function. If f is C 1 , then,
at a point (x, y) U, we can define partial derivatives Dx f (x, y) L(X, Z) and
Dy f (x, y) L(Y, Z) to be the derivatives of the functions
f (, y) : { x X : (x, y) U } Z
at x and y respectively.
Theorem 10.1.2 (Implicit Function Theorem). Suppose that p = n. (That is Y
and Z have the same dimension.) If U X Y is open, f : U Z is C r ,
128
10.2
The first order of business is to show that such partitions of unity exist. The
key idea is the following ingenious construction.
Lemma 10.2.2. There is a C function : R R with (t) = 0 for all t 0 and
(t) > 0 for all t > 0.
Proof. Let
(t) :=
0,
t 0,
1/t
e
, t > 0.
129
Standard facts of elementary calculus can be combined inductively to show that for
each r 1 there is a polynomial Pr such that (r) (t) is Pr (1/t)e1/t if t > 0. Since
the exponential function dominates any polynomial, it follows that (r) (t)/t 0 as
t 0, so that each (r) is differentiable at 0 with (r+1) (0) = 0. Thus is C .
Note that for any open rectangle
x 7
Y
i
Qm
i=1 (ai , bi )
Rm the function
(xi ai )(bi xi )
m
Y
i=1
and Qj,k =
m
Y
i=1
(ki 2)/2j , (ki + 3)/2j .
The cover consists of those Qj,k such that the closure of Qj,k is contained in some
U and, if j > 0, there is no such that the closure of Qj,k is contained in U .
Consider a point x U. The last requirement implies that x has a neighborhood
that intersects only finitely many cubes in the collection, which is to say that the
collection is locally finite.
For any j the Qj,k cover Rm , so there is some k such that x Qj,k , and if j
is sufficiently small, then the closure of Qj,k is contained in some U . If Qj,k is
not in the collection, then the closure of Qj,k is contained in some U . Define k
by letting ki be ki /2 or (ki + 1)/2 according to whether ki is even or odd. Then
Qj,k Qj1,k Qj,k . Repeating this leads eventually to an element of the collection
that contains x, so the collection is indeed a cover of U.
Imposing a coordinate system on X, then combining the observations above,
proves that:
130
Proof. Let
S {Fx : Ux Y }xS be as above. Let { }B be a C partition of unity
for U = x Ux subordinate to {Ux }. For each chooseP
an x such that the closure
of { x : (x) > 0 } is contained in Ux , and let F := Fx : U Y . Then
F is C r because each point in U has a neighborhood in which it is a finite sum of
C r functions. For x S we have
X
X
F (x) =
(x) Fx (x) =
(x) f (x) = f (x).
10.3. MANIFOLDS
10.3
131
Manifolds
132
(ii) Any open subset (including the empty set) of an m-dimensional affine subspace of Rk is an m-dimensional C r manifold. More generally, an open subset
of an m-dimensional C r manifold is itself an m-dimensional C r manifold.
(iii) If U Rm is open and : U Rkm is C r , then the graph
Gr() := { (x, (x)) : x U } Rk
of is an m-dimensional C r manifold, because : x 7 (x, (x)) is a C r
parameterization.
10.4
Smooth Maps
Let M Rk be an m-dimensional C r manifold, and let N R be an ndimensional C r manifold. We have already defined what it means for a function
f : M N is C r to be C r : there is an open W Rk that contains M and a C r
function F : W R such that F |M = f . The following characterization of this
condition is technically useful and conceptually important.
Proposition 10.4.1. For a function f : M N the following are equivalent:
(a) f is C r ;
(b) for each p M there are C r parameterizations : U M and : V N
such that p (U), f ((U)) (V ), and 1 f is a C r function;
(c) 1 f is a C r function whenever : U M and : V N are C r
parameterizations such that f ((U)) (V ).
Proof. Because compositions of C r functions are C r , (a) implies (c), and since each
point in a manifold is contained in the image of a C r parameterization, it is clear
that (c) implies (b). Fix a point p M and C r parameterizations : U M and
: V N with p (U) and f ((U)) (V ). Lemma 10.3.1 implies that 1
and 1 are C r , so ( 1 f ) 1 is C r on its domain of definition. Since
p was arbitrary, we have shown that f is locally C r , and Proposition 10.2.6 implies
that f is C r . Thus (b) implies (a).
There is a more abstract approach to differential topology (which is followed in
Hirsch (1976)) in which an m-dimensional C r manifold is a topological space M
together with a collection { : U M }A , where each is a homeomorphism
betweenS an open subset U of an m-dimensional vector space and an open subset
r
of M, (U ) = M, and for any , A, 1
is C on its domain of
definition. If N with collection { : V : N } is an n-dimensional C r manifold,
a function f : M N is C r by definition if, for all and , 1 f is a C r
function on its domain of definition.
The abstract approach is preferable from a conceptual point of view; for example, we cant see some Rk that contains the physical universe, so our physical
theories should avoid reference to such an Rk if possible. (Sometimes Rk is called
133
the ambient space.) However, in the abstract approach there are certain technical
difficulties that must be overcome just to get acceptable definitions. In addition,
the Whitney embedding theorems (cf. Hirsch (1976)) show that, under assumptions that are satisfied in almost all applications, a manifold satisfying the
abstract definition can be embedded in some Rk , so our approach is not less general
in any important sense. From a technical point of view, the assumed embedding
of M in Rk is extremely useful because it automatically imposing conditions such
as metrizability and thus paracompactness, and it allows certain constructions that
simplify many proofs.
There is a category of C r manifolds and C r maps between them. (This can
be proved from the definitions, or we can just observe that this category can be
obtained from the category of subsets of finite dimensional vector spaces and C r
maps between them by restricting the objects.) The notion of isomorphism for this
category is:
Definition 10.4.2. A function f : M N is a C r -diffeomorphism if f is a
bijection and f and f 1 are both C r . If such an f exists we say that M and N are
C r diffeomorphic.
If M and N are C r diffeomorphic we will, for the most part, regard them as
two different realizations of the same object. In this sense the spirit of the
definition of a C r manifold is that the particular embedding of M in Rk is of no
importance, and k itself is immaterial.
10.5
There are many notions of derivative in mathematics, but invariably the term
refers to a linear approximation of a function that is accurate up to first order.
The first step in defining the derivative of a C r map between manifolds is to specify
the vector spaces that serve as the linear approximations domain and range.
Fix an m-dimensional C r manifold M Rk . Throughout this section, when we
refer to a C r parameterization : U M, it will be understood that U is an open
subset of the m-dimensional vector space X.
Definition 10.5.1. If : U M is a C 1 parameterization and p = (x), then the
tangent space of M at p is the image of this linear transformation D(x) : X
Rk .
We should check that this does not depend on the choice of . If : U M is
a second C 1 parameterization with (x ) = p, then the chain rule gives D (x ) =
D(x)D(1 )(x ), so the image of D (x ) is contained in the image of D(x).
We can combine the tangent spaces at the various points of M:
Definition 10.5.2. The tangent bundle of M is
[
T M :=
{p} Tp M Rk Rk .
pM
134
135
For the categorically minded we mention that Proposition 10.5.4 and the last
three results can be summarized very succinctly by saying that if r 2, then T
is a functor from the category of C r manifolds and C r maps between them to the
category of C r1 manifolds and C r1 maps between them. Again, we will not use
this language later, so in a sense you do not need to know what a functor is, but
categorical concepts and terminology are pervasive in modern mathematics, so it
would certainly be a good idea to learn the basic definitions.
Lets relate the definitions above to more elementary notions of differentiation.
Consider a C 1 function f : (a, b) M and a point t (a, b). Formally Df (t) is
a linear function from Tt (a, b) to Tf (t) M, but thinking about things in this way is
usually rather cumbersome. Of course Tt (a, b) is just a copy of R, and we define
f (t) = Df (t)1 Tf (t) M, where 1 is the element of Tt (A, b) corresponding to 1 R.
When M is an open subset of R we simplify further by treating f (t) as a number
under the identification of Tf (t) M with R. In this way we recover the concept of
the derivative as we first learned it in elementary calculus.
136
10.6
Submanifolds
For almost any kind of mathematical object, we pay special attention to subsets,
or perhaps substructures of other sorts, that share the structural properties of the
object. One only has to imagine a smooth curve on the surface of a sphere to see
that such substructures of manifolds arise naturally. Fix a degree of differentiability
1 r . If M Rk is an m-dimensional C r manifold, N is an n-dimensional
that is also embedded in Rk , and N M, then N is a C r submanifold of M. The
integer m n is called the codimension of N in M.
The reader can certainly imagine a host of examples, so we only mention one
that might easily be overlooked because it is so trivial: any open subset of M
is a C r manifold. Conversely, any codimension zero submanifold of M is just an
open subset. Evidently submanifolds of codimension zero are not in themselves
particularly interesting, but of course they occur frequently.
Submanifolds arise naturally as images of smooth maps, and as solution sets of
systems of equations. We now discuss these two points of view at length, arriving
eventually at an important characterization result. Let M Rk and N R be C r
manifolds that are m- and n-dimensional respectively, and let f : M N be a C r
function. We say that p M is:
(a) an immersion point of f if Df (p) : Tp M Tf (p) N is injective;
(b) a submersion point of f if Df (p) is surjective;
(c) a diffeomorphism point of f is Df (p) is a bijection.
There are now a number of technical results. Collectively their proofs display the
inverse function and the implicit function theorem as the linchpins of the analysis
supporting this subject.
Proposition 10.6.1. If p is an immersion point of f , then there is a neighborhood
V of p such that f (V ) is an m-dimensional C r submanifold of N. In addition
Df (p) : Tp M Tf (p) f (V ) is a linear isomorphism
Proof. Let : U M be a C r parameterization for M whose image contains
p, and let x = 1 (p). The continuity of the derivative implies that there is a
neighborhood U of x such that for all x U the rank of D(f )(x ) is m. Let
X R be the image of Df (p), and let : R X be the orthogonal projection.
Possibly after replacing U with a suitable smaller neighborhood of x, the inverse
function theorem implies that f |U is invertible. Let V = (U ). Now f |U
is an embedding because its inverse is ( f |U )1 . Lemma 10.3.1 implies
that the inverse of f is also C r , so, for every x U the rank of D(f )(x ) is m,
so f (V ) = f ((U )) satisfies Definition 10.3.2.
The final assertion follows from Df (p) being injective while Tp M and Tf (p) (f (V )
are both m-dimensional.
Proposition 10.6.2. If p is a submersion point of f , then there is a neighborhood
U of p such that f 1 (f (p)) U is a (m n)-dimensional C r submanifold of M. In
addition Tp f 1 (q) = ker Df (p).
10.6. SUBMANIFOLDS
137
Proof. Let : U M be a C r parameterization whose image is an open neighborhood of p, let w0 = 1 (p), and let : Z Rn be a C r coordinate chart for an
open neighborhood Z N of f (p). Without loss of generality we may assume that
f ((U)) Z. Since D(w0 ) and D(f (p)) are bijections,
D( f )(w0 ) = D(f (p)) Df (p) D(w0)
is surjective, and the vector space containing U can be decomposed as X Y where
Y is n dimensional and Dy ( f )(w0 ) is nonsingular. Let w0 = (x0 , y0 ). The
implicit function theorem gives an open neighborhood V X containing x0 , an
open W U containing w0 , and a C r function g : V Y such that g(x0 ) = y0 and
{ (x, g(x)) : x V } = { w W : f ((w)) = f (p) }.
Then
{ (x, g(x)) : x V } = f 1 (f (p)) (W )
138
139
10.6. SUBMANIFOLDS
(The third equality follows from the final assertion of Proposition 10.6.2, and the
fourth is the transversality assumption.) Thus p is a submersion point of . Since
p is an arbitrary point of f 1 (P ) the claim follows from Theorem 10.6.5.
We now have
Tp f 1 (P ) = ker D(p) = ker(D(f (p)) Df (p))
= Df (p)1 (ker D(p)) = Df (p)1 (Tf (p) P )
where the first and last equalities are from Proposition 10.6.2.
140
10.7
Tubular Neighborhoods
and : (q, v) 7 q + v
141
()
for all s S with e(f (s), f (s)) s . For each s chooseSan open Us X such that
s Us U(s)/2 (s) and f (Us ) Us /3 (f (s)). Let U = sS Us . We will show that
f |U is injective with continuous inverse.
Consider s, s S and y, y Y with e(f (s), y) < s /3 and e(f (s ), y ) < s /3.
We claim that if y = y , then () holds: otherwise e(f (s), f (s)) > s , s , so that
e(y, y ) e(f (s), f (s)) e(f (s), y) e(f (s ), y )
> ( 21 e(f (s), f (s)) s /3) + ( 21 e(f (s), f (s)) s /3) 16 (s + s ).
142
Proof of the Tubular Neighborhood Theorem. The inverse function theorem implies
that each point (q, 0) in the zero section of N has a neighborhood Nq such that
then is an embedding,
|Ns is a C r1 diffeomorphism. If is in the last result,
S
r1
and its inverse is C
differentiable because U q Nq .
We now develop several applications of the tubular neighborhood theorem. Let
M be an m-dimensional C r -manifold.
Theorem 10.7.6. For any S M, C r1 (S, N) is dense in CS (S, N).
Proof. Proposition 10.2.7 implies that C r1 (S, V ) is dense in CS (S, V ), and Proposition 5.5.3 implies that f 7 1 f is continuous.
Recall that a topological space X is locally path connected if, for each x X,
each neighborhood U of x contains a neighborhood V such that for any x0 , x1 V
there is a continuous path : [0, 1] U with (0) = x0 and (1) = x1 . For an
open subset of a locally convex topological vector space, local path connectedness
is automatic: any neighborhood of a point contains a convex neighborhood.
Theorem 10.7.7. For any S M, CS (S, N) is locally path connected.
Proof. Fix a neighborhood U CS (S, N) of a continuous f : S N. The definition
of the strong topology implies that there is an open W S N such that f { f
C(S, N) : Gr(f ) W } U. Lemma 10.7.4 implies that there is a continuous
: N (0, ) such that U(y) (y) V for all y N and (x, (1 (z))) W for
all x S and z U(f (x)) (f (x)). Let W = { (x, y) W : y U(f (x)) (f (x)) }. For
f0 , f1 C(S, N) with Gr(f0 ), Gr(f1 ) W we define h by setting
h(x, t) = 1 ((1 t)f0 (x) + tf1 (x)) .
If f0 and f1 are C r , so that they are the restrictions to S of C r functions defined
on open supersets of S, then this formula defines a C r extension of h to an open
superset of S [0, 1], so that h is C r .
143
It is easy to see (and not hard to compute formally using the chain rule) that
D
(q, 0) = IdTq N IdTq N under the natural identification of T(q,0) (T N) with Tq N
Tq N. The inverse function theorem implies that after replacing Nq with a smaller
neighborhood of (q, 0), the restriction of
to Nq is a diffeomorphism onto its image.
We can now proceed as in the proof of the tubular neighborhood theorem.
The following construction simulates convex combination.
Proposition 10.7.9. There is a neighborhood W of the diagonal in N N and a
continuous function c : W [0, 1] N such that:
(a) c((q, q ), 0) = q for all (q, q ) W ;
(b) c((q, q ), 1) = q for all (q, q ) W ;
(c) c((q, q), t) = q for all q N and all t.
Proof. The tubular neighborhood gives an open neighborhood U of the zero section
in N such that if : N Rk is the map (q, v) = q + v, then |U is a homeomorphism between U and (U). Let : N N be the projection for the normal
bundle. Let
W = { (q, q ) N N : (1 t)q + tq (U) for all 0 t 1 },
and for (q, q ) W and 0 t 1 let
c((q, q ), t) = (|U )1 ((1 t)q + tq ) .
10.8
144
145
)
0 there is some i such that (h
(
x) 6= 0; after reindexing we may assume that
xi
m
R be the function
i = m. Let : W
(Of
course there is no difficulty showing that D (x)
If r 2, then T is an
injective C r1 -immersion whose image is open in T M, so it is a C r -embedding.
Since T M is covered by the images of maps such as T , it is indeed a C r1 manifold.
146
10.9
In order to study the behavior of fixed points under homotopy, we will need
to understand the structure of h1 (q) when M and N are manifolds of the same
dimension,
h : M [0, 1] N
1
If : V N is a C r parameterization for N whose image contains f (W ), then 1 has a C r
extension, because that is what it means for a function on a possibly nonopen domain to be C r ,
and this extension can be composed with to give f.
147
(s a )
(s a ) + (b s)
for sufficiently small , > 0 gives a C function : (a, b) [0, 1] with (s) = 0
for all s (a, a + ), (s) = 1 for all s (b , b), and (s) > 0 for all s such that
0 < (s) < 1. For any real number Q we can define : (a, b) R by setting
(s) = (1 (s))(s Q) + (s)f (s).
Clearly this will be satisfactory if (s) > 0 for all s. A brief calculation gives
(s) = 1 + (s)(f (s) 1) + (s)(Q + f (s) s)
= (1 (s))(1 f (s)) + f (s) + (s)(Q + f (s) s).
If Q is larger than the upper bound for s f (s), then (s) > 0 when (s) is close
to 0 or 1. Since those s for which this is not the case will be contained in a compact
interval on which positive and continuous, hence bounded below by a positive
constant, if Q is sufficiently large then (s) > 0 for all s.
Proof of Proposition 10.9.1. Let M be a nonempty compact connected 1-dimensional
C r manifold. We can pass from a C r atlas for M to a C r atlas whose elements all
have connected domains by taking the restrictions of each element of the atlas to
the connected components of its domain. To be concrete, we will assume that the
domains of the parameterizations are connected subsets of R, i.e., open intervals.
Since we can pass from a parameterization with unbounded domain to a countable
collection of restrictions to bounded domains, we may assume that all domains are
bounded. Since M is compact, any atlas has a finite subset that is also an atlas.
We now have an atlas of the form
{ 1 : (a1 , b1 ) M, . . . , K : (aK , bK ) M }.
148
1 (s + Q), s b1 Q,
(s) = 1 (1 (s)), b1 Q < s < t2 ,
2 (s),
s t2 .
We have 1 (s) = s + Q for all s in a neighborhood of b1 Q and (s) = 2 (s)
for all s close to t2 . Therefore is a C r function. Each point in its domain
has a neighborhood such that the restriction of to that neighborhood is a C r
parameterization for M, which implies that if maps open sets to open sets. If it
was injective, it would be a C r coordinate chart whose image was the union of the
images of 1 and 2 , which would contradict the minimality of K.
Therefore is not injective. Since 1 and 2 are injective, there must be s <
b1 c such that (s) = (s ) for some s > t1 . Let s0 be the supremum of such
s. If (s0 ) = (s ) for some s > t1 , then the restrictions of to neighborhoods of
s0 and s would both map diffeomorphically onto some neighborhood of this point,
which would give a contradiction of the definition of s0 . Therefore (s0 ) is in the
closure of (((t1 , b2 )), but is not an element of this set, so it must be lims b2 (s ).
Arguments similar to those given above imply that there are , > 0 such that the
images of |(b2 ,b2 ) and |(s0 ,s0 ) are the same, and the C r diffeomorphism
g = (|(s0 ,s0 ) )1 |(b2 ,b2 )
is increasing. Applying the lemma above again, there is a real number R and an
increasing C r diffeomorphism : (b2 , b2 ) (b2 R, s0 ) such that (s) = sR
for s near b2 and (s) = g(s) for s near b2 .
We now define : [s0 , s0 + R) M by setting
(
(s),
s0 s b2 ,
(s) =
1
( (s R)), b2 < s < s0 + R.
Then agrees with near b2 , so it is C r , and it agrees with (s R) near
s0 + R, so it can be construed as a C r function from the circle (thought of R modulo
149
R) to M. This function is easily seen to be injective, and it maps open sets to open
sets, so its image is open, but also compact, hence closed. Since M is connected,
its image must be all of M, so we have constructed te desired C r diffeomorphism
between the circle and M.
The argument for a compact connected one dimensional C r -manifold with
nonempty boundary is similar, but somewhat simpler, so we leave it to the reader.
Although it will not figure in the work here, the reader should certainly be
aware that the analogous issues for higher dimensions are extremely important in
topology, and mathematical culture more generally. In general, a classification of
some type of mathematical object is a description of all the isomorphism classes
(for whatever is the appropriate notion of isomorphism) of the object in question.
The result above classifies compact connected 1-dimensional C r manifolds.
The problem of classifying oriented surfaces (2-dimensional manifolds) was first
considered in a paper of Mobius in 1870. The classification of all compact connected
surfaces was correctly stated by van Dyke in 1888. This result was proved for
surfaces that can be triangulated by Dehn and Heegaard in 1907, and in 1925 Rado
showed that any surface can be triangulated.
After some missteps, Poincare formulated a fundamental problem for the the
classification of 3-manifolds: is a simply connected compact 3-manifold necessarily
homeomorphic to S 3 ? (A topological space X is simply connected if it is connected and any continuous function f : S 1 = { (x, y) R2 : x2 + y 2 = 1 } X has
a continuous extension F : D 2 = { (x, y) R2 : x2 + y 2 1 } X.) Although
Poincare did not express a strong view, this became known as the Poincar
e conth
jecture, and over the course of the 20 century, as it resisted solution and the four
color theorem and Fermats last theorem were proved, it became perhaps the most
famous open problem in mathematics. Curiously, the analogous theorems for higher
dimensions were proved first, by Smale in 1961 for dimensions five and higher, and
by Freedman in 1982 for dimension four. Finally in late 2002 and 2003 Perelman
posted three papers on the internet that sketched a proof of the original conjecture.
Over the next three years three different teams of two mathematicians set about
filling in the details of the argument. In the middle of 2006 each of the teams posted
a (book length) paper giving a complete argument. Although Perelmans papers
were quite terse, and many details needed to be filled in, all three teams agreed
that all gaps in his argument were minor.
Chapter 11
Sards Theorem
The results concerning existence and uniqueness of systems of linear equations
have been well established for a long time, of course. In the late 19th century Walras recognized that the system describing economic equilibria had (after recognizing
the redundant equation now known as Walras law) the same number of equations
and free variables, which suggested that typically economic equilibria should be
isolated and also robust, in the sense that the endogenous variables will vary continuously with the underlying parameters in some neighborhood of the initial point. It
was several decades before methods for making these ideas precise were established
in mathematics, and then several more decades elapsed before they were imported
into theoretical economics.
The original versions of what is now known as Sards theorem appeared during
the 1930s. There followed a process of evolution, both in the generality of the
result and in the method of proof, that culminated in the version due to Federer
(see Section 11.3.) Our treatment here is primarily based on Milnor (1965), fleshed
out with some arguments from Sternberg (1983), which (in its first edition) seems
to have been Milnors primary source. While not completely general, this version of
the result is adequate for all of the applications in economic theory to date, many
of which are extremely important.
Suppose 1 r , and let f : U Rn be a C r function, where U Rm
is open. If f (x) = y and Df (x) has rank n, then the implicit function theorem
(Theorem 10.1.2) implies that, in a neighborhood of x, f 1 (y) can be thought of
as the graph of a C r function. Intuition developed by looking at low dimensional
examples suggests that for typical values of y this pleasant situation will prevail
at all elements of f 1 (y), but even in the case m = n = 1 one can see that there can
be a countable infinity of exceptional y. Thus the difficulty in formulating this idea
precisely is that we need a suitable notion of a small subset of Rn . This problem
was solved by the theory of Lesbesgue measure, which explains the relatively late
date at which the result first appeared.
Measure theory has rather complex foundations, so it preferable that it not be
a prerequisite. Thus it is fortunate that only the notion of a set of measure zero
is required. Section 11.1 defines this notion and establishes its basic properties.
One of the most important results in measure theory is Fubinis theorem, which,
roughly speaking, allows functions to be integrated one variable at a time. Section
150
151
11.2 develops a Fubini-like result for sets of measure zero. With these elements
in place, it becomes possible to state and prove Sards theorem in Section 11.3.
Section 11.4 explains how to extend the result to maps between sufficiently smooth
manifolds.
The application of Sards theorem that is most important in the larger scheme of
this book is given in Section 11.5. The overall idea is to show that any map between
manifolds can be approximated by one that is transversal to a given submanifold
of the range.
11.1
For each n there is a positive constant such that the volume of a ball in Rn is
that constant times r n , where r is the radius of the ball. Without knowing very
much about the constant, we can still say that sets satisfying the following definition
are small.
Definition 11.1.1. A set S Rm has measure zero if, for any > 0, there is
k
a sequence {(xj , rj )}
j=1 in R (0, 1) such that
S
[
j
rjm < .
2j with
(2j )m < , then there is a covering by balls of radius mj with
j
Qm
P
m
m
i=1 [ai , bi ] because we can
j ( mj ) < ( m/2) . We can also use rectangles
cover such a rectangle with a collection of cubes of almost the same total volume;
from the point of view of our methodology it is important to recognize that we
know this as a fact of arithmetic (and in particular the distributive law) rather
than as prior knowledge concerning volume.
The rest of this section develops a few basic facts. The following property of
sets of measure zero occurs frequently in proofs.
Lemma 11.1.2. If S1 , S2 , . . . Rm are sets of measure zero, then S1 S2 . . . has
measure zero.
Proof. For given take the union of a countable cover of S1 by rectangles of total
volume < /2, a countable cover of S2 by rectangles of total volume < /4, etc.
It is intuitively obvious that a set of measure zero cannot have a nonempty
interior, but our methodology requires that we forget everything we know about
volume, using only arithmetic to prove it.
Lemma 11.1.3. If S has measure zero, its interior is empty, so its complement is
dense.
152
Proof. Suppose that, on the contrary, S has a nonempty interior. Then it contains a
closed cube C, say of side length
P 2. Fixing > 0, suppose that S has a covering by
cubes of side length 2j with j (2j )m < . Then it has a covering by open cubes
Cj of side length 3j , and there is a finite subcover of C. For some large integer K,
Q
ij ij +1
consider all standard cubes of the form m
j=1 [ K , K ]. For each cube in our finite
subcover, let Dj be the union of all such standard cubes contained in Cj , and let nj
be the number of such cubes. Let D be the union of all standard cubes containing a
point in C, and let n be the number of them. Simply as a matter of counting (that
is to say, without reference to any theory of volume)Swe have nj /K m P
(3j )m and
n/K m (2)m . If K is sufficiently large, then D j Dj , so that n j nj and
(2)m n/K m
X
j
nj /K m
X
j
(3j )m (3/2)m ,
If {(xj , rj )}
j=1 is a sequence such that
SC
rjm < ,
(Brj )m < B m .
then
f (S C)
[
j
11.2
153
154
Proof. Fix > 0, and choose < /2(b1 a1 ). Since P (C) P (C), it has one
dimensional measure zero, and since it is closed, hence compact, it can be covered by
the union J of finitely many open intervals of total length /2(b2 a2 ) (bm am ).
In this way { x C : x1 J } is covered by a union of open rectangles of total
volume /2.
For each t
/ J we can choose a finite union of rectangles in Rm1 of total volume
less than that covers C(t), and these will also cover C(t ) for all t in some open
interval around t. Since [a1 , b1 ] \ J is compact, it is covered by a finite collection of
such intervals, and it is evident that we can construct a cover of { x C : x1
/J}
of total volume less than /2.
Lemma 11.2.5. If C has measure zero, then P (C) has measure zero.
S
Proof. Since P (C) = n=1,2,... P1/n (C), it suffices to show that P (C) has measure
zero for any > 0. For any > 0 there is a covering of C by finitely many rectangles
of total volume less than . For each t there is an induced covering C(t) be a finite
collection of rectangles, and there is an induced covering of [a1 , b1 ]. The total length
of intervals with induced coverings of total volume greater than cannot exceed
/.
11.3
Sards Theorem
We now come to this chapters central result. Recall that a critical point of a
C function is a point in the domain at which the rank of the derivative is less than
the dimension of the range, and a critical value is a point in the range that is the
image of a critical point.
1
155
closed set, so it is covered by countably many compact sets, each of which is covered
by finitely many such neighborhoods, and consequently it has a countable cover by
such neighborhoods.
f1
After reindexing we may assume that x
(x) 6= 0. Let V be a neighborhood of
1
f1
x in which x1 does not vanish. Let h : V Rm be the function
h(x) := (f1 (x), x2 , . . . , xm ).
The matrix of partial derivatives of h at x is
f
f1
1
(x) x
(x)
x1
2
0
1
..
.
..
.
0
f1
(x)
xm
0
..
.
so the inverse function theorem implies that, after replacing V with a smaller neighborhood of x, h is a diffeomorphism onto its image. The chain rule implies that
the critical values of f are the critical values of g = f h1 , so we can replace f
with g, and g has the additional property that g1 (z) = z1 for all z in its domain.
The upshot of this argument is that we may assume without loss of generality that
f1 (x) = x1 for all x V .
For each t R let V t := { w Rm1 : (t, w) V }, let f t : V t Rn1 be the
function
f t (w) := (f2 (t, w), . . . , fn (t, w)),
and let C t be the set of critical points of f t . The matrix of partial derivatives of f
at x V is
1
0
0
f2 (x) f2 (x) f2 (x)
x2
xm
x1
..
..
.. ,
.
.
.
fn
(x)
x1
fn
(x)
x2
fn
(x)
xm
Since the result is known to be true with (m, n) replaced by (m 1, n 1), each
f t (C t ) has (n 1)-dimensional measure zero. In addition, the continuity of the
relevant partial derivatives implies that C \ C1 is locally closed, so Proposition
11.2.1 implies that f (C V ) has measure zero.
i+1
f
(x), x2 , . . . , xm ).
h(x) := ( xs x
s
1
156
After reindexing we may assume that si+1 = 1, so that the matrix of partial derivatives of h at x is triangular with nonzero diagonal entries. By the inverse function
theorem the restriction of h to some neighborhood V of x is a C diffeomorphism.
Let g := f (h|V )1 . Then h(V Ci ) {0} Rm1 . Let
g0 : { y Rm1 : (0, y) h(V ) } Rn
be the map g0 (y) = g(0, y). Then f (V (Ci \Ci+1 )) is contained in the set of critical
values of g0 , and the latter set has measure zero because the result is already known
when (m, n) is replaced by (m 1, n).
Proof of (c): Since U can be covered by countably many compact cubes, it suffices
to show that f (Cr I) has measure zero whenever I U is a compact cube. Since I
is compact and the partials of f of order r are continuous, Taylors theorem implies
that for every > 0 there is > 0 such that
kf (x + h) f (x)k khkr
whenever x, x + h I with x Cr and khk < . Let L be the side length of I. For
each integer d > 0divide I into dm subcubes of side length L/d. The diameter of
such a subcube is mL/d. If this quantity is less than and the subcube
contains a
r n n
2( mL) , and may be arbitrarily small.
157
Note that this definition makes perfect sense even if is not an integer! Let
U Rm be open, and let f : U Rn be a C r function. For 0 p < m let Rp
be the set of points x M such that the rank of Df (x) is less than or equal to
p. The most general and sophisticated version of Sards theorem, due to Federer,
states that f (Rp ) has -dimensional measure zero for all > p + mp
. A beautiful
r
informal introduction to the circle of ideas surrounding these concepts, which is the
branch of analysis called geometric measure theory, is given by Morgan (1988). The
proof itself is in Section 3.4 of Federer (1969). This reference also gives a complete
set of counterexamples showing this result to be best possible.
11.4
158
11.5
Genericity of Transversality
159
160
(a) Gr(f) A;
(b) f|M \Z = f |M \Z ;
(c) f|W Z is smooth;
(d) f CK P .
Proof. Let
P = 1 (P ),
g = 1 f ,
C = 1 (C),
= 1 (W ),
W
= 1 (K),
K
Z = 1 (Z).
Let Z and g be the set and function whose existence is guaranteed by the last result.
Set Z = (Z ), and define f by specifying that f agrees with f on M \ (U), and
f|(U ) = g 1 .
Clearly f has all the desired properties.
In order to apply this we need to have an ample supply of smooth parameterizations for N that are aligned with h.
Lemma 11.5.6. Each point q f (M) is contained in the image of a smooth
parameterization that is aligned with h.
Proof. Let : V N be any smooth parameterization whose image contains q,
and let y = 1 (q). Let : W S be a smooth parameterization whose image
contains h(q); we can replace V with a smaller open set containing y, so we may
V )) (W ).
assume that h((
Since h f is a submersion, the rank of Dh(q) is s, and consequently the rank
of D( 1 h )(y)
is also s. After some reindexing, y is a regular point of
))), y , . . . , y ).
: y 7 ( 1 (h((y
n
s+1
Applying the inverse function theorem, a smooth parameterization whose image
contains q that is aligned with h is given by letting be the restriction of 1
to some neighborhood of ( 1 (h(q)), ys+1, . . . , yn ).
Proof of Theorem 11.5.3. Any open subset of M is a union of open subsets whose
closures are compact. In view of this fact and the last result, M is covered by the
sets (U) where : U M is a smooth parameterization with f ((U)) (V ) for
some smooth parameterization : V N that is aligned with h, and the closure
of (U) is compact. Since M is paracompact, there is a locally finite cover by the
images of such parameterizations, and since M is separable, this cover is countable.
That is, there is a sequence 1 : U1 M, 2 : U2 M, . . . whose images cover
M, such that for each i, the closure of i (Ui ) is compact and there is a smooth
parameterization i : Vi N that is aligned with h such that f (i (Ui )) i (Vi ).
We claim that there is a sequence K1 , K2 , . . . of compact subsets of M that cover
M, with Ki i (Ui ) for each i. For p M let (p) be the maximum such that
161
U (p) i (Ui ) for some i, and let ip be an integer that attains the maximum. Then
: M (0, ) is a continuous function. For each i let i := minpi (Ui ) (p), and
let
Ki = { p i (Ui ) : Ui (p) i (Ui ) }.
is unchanging after some point. Thus the sequence f1 , f2 , . . . has a well defined limit
that is smooth and transversal to P , and whose graph is contained in A.
We now turn to the proof of Proposition 11.5.4. The main idea is to select a
suitable member from a family of perturbations of g. The following lemma isolates
the step in the argument that uses Sards theorem.
Lemma 11.5.7. If U Rm and B Rns are open, P is a p-dimensional smooth
submanifold of Rn , and G : U B Rn is smooth and transversal to P , then for
almost every b B the functions gb = G(, b) : U N is transversal to P .
Proof. Let Q = G1 (P ). By the transversality theorem, Q is a smooth manifold,
of dimension (m + (n s)) (n p) = m + p s. Let be the natural projection
U B B. Sards theorem implies that almost every b B is a regular value of
|Q . Fix such a b. We will show that gb is transversal to P .
Fix x gb1 (P ), set q = gb (x), and choose some y Tq N. Since G is transversal
to P there is a u T(x,b) (U B) such that y is the sum of DG(x, b)u and an element
of Tq P . Let u = (v, w) where v Rm and w Rns . Since (x, b) is a regular point
of |Q , there is a u T(x,b) Q such that D|Q (x, b)u = w. Let u = (v , w).
Then Tq P contains DG(x, b)u , so it contains
DG(x, b)u y + DG(x, b)u = DG(x, b)(v + v , 0) y = Dgb(x)(v + v ) y.
Thus y is the sum of Dgb (x)(v + v ) and an element of Tq P , as desired.
Proof of Proposition 11.5.4. For x U let (x)
162
For any (x, b) the image of DG(x, b) contains the image of Dg(x), so, since g|Y P ,
we have G Y B P . Since g s is a submersion, at every (x, b) such that (x) > 0 the
image of DG(x, b) is all of Rn , so G (U \Y )B P . Therefore G P . The last result
implies that for some b B, gb = G(, b) is transversal to P . Evidently gb agrees
with g on Y .
Chapter 12
Degree Theory
Orientation is an intuitively familiar phenomenon, modelling, among other things,
the fact that there is no way to turn a left shoe into a right shoe by rotating it, but
the mirror image of a left shoe is a right shoe. Consider that when you look at a
mirror there is a coordinate system in which the map taking each point to its mirror
image is the linear transformation (x1 , x2 , x3 ) 7 (x1 , x2 , x3 ). It turns out that the
critical feature of this transformation is that its determinant is negative. Section
12.1 describes the formalism used to impose an orientation on a vector space and
consistently on the tangent spaces of the points of a manifold, when this is possible.
Section 12.2 discusses two senses in which an orientation on a given object
induces a derived orientation: a) an orientation on a -manifold induces an orientation of its boundary; b) given a smooth map between two manifolds of the same
dimension, an orientation of the tangent space of a regular point in the domain
induces an orientation of the tangent space of that points image. If both manifolds
are oriented, we can define a sense in which the map is orientation preserving or
orientation reversing by comparing the induced orientation of the tangent space of
the image point with its given orientation.
In Section 12.3 we first define the smooth degree of a smooth (where smooth
now means C ) map over a regular value in the range to be the number of preimages
of the point at which the map is orientation preserving minus the number of points at
which it is orientation reversing. Although the degree for smooth functions provides
the correct geometric intuition, it is insufficiently general. The desired generalization is achieved by approximating a continuous function with smooth functions,
and showing that any two sufficiently accurate approximations are homotopic, so
that such approximations can be used to define the degree of the given continuous
function. However, instead of working directly with such a definition, it turns out
that an axiomatic characterization is more useful.
12.1
Orientation
The intuition underlying orientation is simple enough, but the formalism is a bit
heavy, with the main definitions expressed as equivalence classes of an equivalence
relation. We assume prior familiarity with the main facts about determinants of
matrices.
163
164
No doubt most readers are well aware that a linear automorphism (that is, a
linear transformation from a vector space to itself) has a determinant. What we
mean by this is that the determinant of the matrix representing the transformation
does not depend on the choice of coordinate system. Concretely, if L and L are the
matrices of the transformation in two coordinate systems, then there is a matrix U
(expressing the change of coordinates) such that L = U 1 LU, so that
|L | = |U 1 LU| = |U|1 |L| |U| = |L|.
Let X be an m-dimensional vector space. An ordered basis of X is an ordered
m-tuple (v1 , . . . , vm ) of linearly independent vectors in X. Mostly we will omit
the parentheses, writing v1 , . . . , vm when the interpretation is clear. If v1 , . . . , vm
and v1 , . . . , vm
are ordered bases, we say that they have the same orientation if
v1 , . . . , vm and v1 , . . . , vm
must have the same orientation, so there are precisely
two equivalence classes. An orientation for X is one of these equivalence classes.
An oriented vector space is a vector space for which one of the two orientations
has been specified. An ordered basis of an oriented vector space is said to be
positively oriented (negatively oriented) if it is (not) an element of the specified
orientation.
Since the determinant is continuous, each orientation is an open subset of the
set of ordered bases of X. The two orientations are disjoint, and their union is the
entire set of ordered bases, so each path component of the space of ordered bases
is contained in one of the two orientations. If the space of ordered bases had more
than two path components, it would be possible to develop an invariant that was
more sophisticated than orientation. But this is not the case.
Proposition 12.1.1. Each orientation of X is path connected.
Proof. Fix a standard basis e1 , . . . , em and some ordered basis v1 , . . . , vm . We
will show that there is a path in the space of ordered bases from v1 , . . . , vm to
either e1 , e1 , . . . , em or e1 , e1 , . . . , em . Thus the space of ordered bases has at
most two path components, and since each orientation is a nonempty union of path
components, each must be a path component.
If i 6= j, then for any t R the determinant of the linear transformation
taking v1 , . . . , vm to v1 , . . . , vi + tvj , . . . , vm is one, so varying t gives a continuous
path in the space of ordered bases. Combining
P such paths, we can find a path
from v1 , . . . , vm to w1 , . . . , wm where wi = j bij ej with bij 6= 0 for all i and j.
Beginning at w1 , . . . , wm , such paths can be combined to eliminate all off diagonal
coefficients, arriving at an ordered basis of the form c1 e1 , . . . , cm em . From here we
12.1. ORIENTATION
165
166
i i ui + uh+1 .
The general result is obtained by applying this in the context of finite collection
of parameterizations that cover .
Proof of Proposition 12.1.3. There are a = t0 < t1 < < tJ1 < tJ = b such
that for each j = 1, . . . , J, the image of |[tj1 ,tj ] is contained in the image of a
smooth parameterization. We may assume that J = 1 because the general case
can obviously be obtained from J applications of this special case. Thus there is
a smooth parameterization : U M whose image contains the image of . Let
= 1 , let := , let uh+1 = D((a))vh+1 (a), and define a moving h-frame
u along by setting
u1 (t) := D((t))v1 (t), . . . , uh (t) := D((t))vh (t).
The last result gives a uh+1 : [a, b] Rm such that uh+1 (a) = uh+1 and (u1 , . . . , uh , uh+1 )
is a moving (h + 1)-frame along . We define vh+1 to [a, b] by setting
vh+1 (t) = D((t))uh+1(t).
167
12.1. ORIENTATION
ij (t)vj (t).
168
12.2
Induced Orientation
An orientation on a manifold induces an orientation on an open subset, obviously. More interesting is the orientation induced on M by an orientation on the
a -manifold M. We are also interested in how an orientation on a point in the
domain of a smooth map between manifolds of equal dimension induces an orientation on the tangent space of the image point. As we will see, this generalizes to the
image point being in an oriented submanifold whose codimension is the dimension
of the domain.
As before we work with an m-dimensional smooth -manifold M with a given orientation. Consider a point p M and a basis v1 , . . . , vm of Tp M with v2 , . . . , vm
Tp M. Of course v2 , . . . , vm is a basis of Tp M. There is a visually obvious sense
in which v1 is either inward pointing or outward pointing that is made precise
by using a parameterization : U M (where U H is open) to determine
whether the first component of D1 (p)v1 is positive or negative. Note that the
sets of inward and outward point vectors are both convex. Our convention will be
that an orientation of Tp M induces an orientation of Tp M according to the rule
that if v1 , . . . , vm is positively oriented and v1 is outward pointing, then v1 , . . . , vm1
is positively oriented.
Does our definition of the induced orientation make sense? There are two issues
to address.
169
define the induced orientation of Tp M, then they give the same induced orientation.
Suppose that v1 and v1 are both outward pointing. Since the set of outward pointing
vectors is convex,
t 7 (1 t)v1 + tv1 , v2 , . . . , vm
taking v1 , v2 , . . . , vm to v1 , v2 . . . , vm
(concretely, vi =
j aij vj ) is (1, 0, . . . , 0),
so the determinant of A is the same as the determinant of its lower right hand
(m 1) (m 1) submatrix, which is the matrix of the linear transformation
taking v2 , . . . , vm to v2 , . . . , vm
. Therefore v1 , . . . , vm has the same orientation as
v2 , . . . , vm
.
We also need to check that what we have defined as the induced orientation of
M is, in fact, an orientation. Consider a path : [a, b] M. Corollary 12.1.5
gives a moving frame (v2 , . . . , vm ) for M along , and Proposition 12.1.3 implies
that it extends to a moving frame (v1 , . . . , vm ) for M along . Suppose that v1 (a) is
outward pointing. By continuity, it must be the case that for all t, v1 (t) is outward
pointing. If we assume that v1 (a), . . . , vm (a) is positively oriented, for the given
orientation, then v2 (a), . . . , vm (a) is positively oriented, for the induced orientation.
In addition, v1 (b), . . . , vm (b) is positively oriented, for the given orientation, so, as
desired, v2 (b), . . . , vm (b) is positively oriented, both with respect to the induced
orientation and with respect to the orientation induced by and v2 (a), . . . , vm (a).
Now suppose that M and N are two m-dimensional oriented smooth manifolds,
now without boundary, and that f : M N is a smooth function. If p is a
regular point of f , we say that f is orientation preserving at p if Df (p) maps
positively oriented bases of Tp M to positively oriented bases of Tf (p) N; otherwise f is
170
v1 (b), v2 , . . . , vm+1
of T(b) M, so we may assume that v2 (b), . . . , vm+1 (b) T(b) M.
Then v1 (a), . . . , vm+1 (a) is a positively oriented basis of T(a) M if and only if
v1 (b), . . . , vm+1 (b) is a positively oriented basis of T(b) M. Since v1 (a) is inward
pointing and v1 (b) is outward pointing, v2 (a), . . . , vm+1 (a) is a positively oriented
basis of T(a) M if and only if v2 (b), . . . , vm+1 (b) is a negatively oriented basis of
T(b) M.
Proposition 12.1.3 implies that there is a moving frame w1 , . . . , wnm along
f : [a, b] P . As we have defined orientation, w1 (a), . . . , wnm (a) is a positively
oriented basis of Tf ((a)) P if and only if w1 (b), . . . , wnm (b) is a positively oriented
basis of Tf ((b)) P , and
171
12.3
The Degree
1
(1) deg
(q) is a singleton {p}
q (f ) = 1 for all (f, q) D (M, N) such that f
and f is orientation preserving at p.
Pr
(2) deg
q (f ) =
i=1 degq (f |Ci ) whenever (f, q) D (M, N), the domain of f
is C, and C1 , . . . , Cr are pairwise disjoint compact subsets of C such that
f 1 (q) C1 . . . Cr \ (C1 . . . Cr ).
(3) deg
q (h0 ) = deg q (h1 ) whenever C M is compact and the homotopy h :
C [0, 1] N is smoothly degree admissible over q.
1
Concretely, deg
(q) at which f is orientation preservq (f ) is the number of p f
1
ing minus the number of p f (q) at which f is orientation reversing.
172
Proof. For (f, q) D(M, N) the inverse function theorem implies that each p
f 1 (q) has a neighborhood that contains no other element of f 1 (q), and since U is
1
compact it follows that f 1 (q) is finite. Let deg
(q)
q (f ) be the number of p f
1
at which f is orientation preserving minus the number of p f (q) at which f is
orientation reversing.
Clearly deg satisfies (1) and (2). Suppose that h : C [0, 1] N is
smoothly degree admissible over q. Let V be a neighborhood of q such that for all
q V :
(a) h1 (q ) U [0, 1];
(c) deg
q (h0 ) = deg q (h0 ) and deg q (h1 ) = deg q (h1 ).
deg
q (f ) = 2 degq (f ) + 1.
In preparation for the next result we show that deg is continuous in a rather
strong sense.
Proposition 12.3.4. If C M is compact, f : C N is continuous, and q
N \ f (C), then are neighborhoods Z C(C, N) of f and V N \ f (C) of q such
that
deg
q (f ) = deg q (f )
173
deg
q (f ) = degq (f ) whenever f Z C (C, N) and q , q V are regular values
of f . Let j : N [0, 1] N be a smooth function with the following properties:
(a) j0 = IdN ;
(b) each jt is a smooth diffeomorphism;
(c) j(y, t) = y for all y N \ V and all t;
(d) j1 (q ) = q .
(Construction of such a j, using the techniques of Section 10.2, is left as an exercise.)
Clearly jt (q ) is a regular value of jt f for all t, so the concrete characterization of
174
Proof. We claim that if deg : D(M, N) Z satisfies (D1)-(D3), then its restriction
to D (M, N) satisfies (1)-(3). For (1) and (2) this is automatic. Suppose
that C M is compact and h : U [0, 1] N is a smoothly degree admissible
homotopy over q. Such a homotopy may be regarded as a continuous function from
[0, 1] to C(U , N). Therefore (D3) implies that degq (ht ) is a locally constant function
of t, and since [0, 1] is connected, it must be constant. Thus (3) holds.
Proposition 11.5.1 implies that for any (f, q) D(M, N) the set of smooth
f : M N that have q as a regular value is dense at f . In conjunction with
Proposition 12.3.4, this implies that the only possibility consistent with (D3) is to
12.4
In Chapter 5 we emphasized restriction to a subdomain, composition, and cartesian products, as the basic set theoretic methods for constructing new functions from
ones that are given. The bahevior of the degree under restriction to a subdomain is
already expressed by (D3), and in this section we study the behavior of the degree
under composition and products. In both cases the result is given by multiplication,
reflecting basic properties of the determinant.
Proposition 12.4.1. If M, N, and P are oriented m-dimensional smooth manifolds, C M and D N are compact, f : C N and g : D P are continuous,
g is degree admissible over r P , and g 1 (r) is contained in one of the connected
components of N \ f (C), then for any q g 1(r) we have
degr (g f ) = degq (f ) degr (g).
Proof. Since C (C, N) and C (D, P ) are dense in C(C, N) and C(D, P ) (Theorem
10.7.6) and composition is a continuous operation (Proposition 5.3.6) the continuity
175
property (D3) of the degree implies that is suffices to prove the claim when f and
g are smooth. Sards theorem implies that there are points r arbitrarily near r
that are regular values of both g and g f , and Proposition 12.3.4 implies that the
relevant degrees are unaffected if r is replaced by such a point, so we may assume
that r has these regularity properties.
For q g 1 (r) let sg (q) be 1 or 1 according to whether g is orientation preserving or orientation reversing at q. For p (g f )1 (q) define sf (p) and sgf (p)
similarly. In view of the chain rule and the definition of orientation preservation
and reversal, sgf (p) = sg (f (p))sf (p). Therefore
X
X
X
deg(g f ) =
sg (f (p))sf (p) =
sg (q)
sf (p)
p(gf )1 (r)
qg 1 (r)
pg 1 (q)
sg (q) degq (f ).
qg 1 (r)
Proof. For reasons explained in other proofs above, we may assume that f and f
are smooth and that q and q are regular values of f and f . For p f 1 (r) let
sf (p) be 1 or 1 according to whether f is orientation preserving or orientation
reversing at p, and define sf (p ) for p f 1 (q ) similarly. Since the determinant
of a block diagonal matrix is the product of the determinants of the blocks, f f
is orientation preserving or orientation reversing at (p, p ) according to whether
sp (f )sp (f ) is positive or negative, so
X
sp (f )sp (f )
deg(q,q ) (f f ) =
(p,p )(f f )1 (q,q )
pf 1 (q)
sp (f )
p f 1 (q )
sp (f ) = degq (f ) degq (f ).
Chapter 13
The Fixed Point Index
We now take up the theory of the fixed point index. For continuous functions defined on compact subsets of Euclidean spaces it is no more than a different rendering
of the theory of the degree; this perspective is developed in Section 13.1.
But we will see that it extends to a much higher level of generality, because the
domain and the range of the function or correspondence have the same topology.
Concretely, there is a property called Commutativity that relates the indices of the
and g : C X are continuous,
two compositions g g and g g where g : C X
and other natural restrictions on this data (that will give rise to a quite cumbersome
definition) are satisfied. This property requires that we extend our framework to
allow comparison across spaces. Section 13.2 introduces the necessary abstractions
and verifies that Commutativity is indeed satisfied in the smooth case. It turns out
that this boils down to a fact of linear algebra that came as a surprise when this
theory was developed.
When we extended the degree from smooth to continuous functions, we showed
that continuous functions could be approximated by smooth ones, and that this gave
a definition of the degree for continuous functions that made sense and was uniquely
characterized by certain axioms. In somewhat the same way Commutativity will
be used, in Section 13.4, to extend from Euclidean spaces to separable ANRs, as
per the ideas developed in Section 7.6. The argument is lengthy, technically dense,
and in several ways the culmination of our work to this point.
The Continuity axiom is then used in Section 13.5 to extend the index to contractible valued correspondences. The underlying idea is the one used to extend
from smooth to continuous functions: approximate and show that the resulting definition is consistent and satisfies all properties. Again, there are many verifications,
and the argument is rather dense.
Multiplication is an additional property of the index that describe its behavior in connection with cartesian products. For continuous functions on subsets of
Euclidean spaces it is a direct consequence of Proposition 12.4.3. At higher levels
of generality it is, in principle, a consequence of the axioms, because those axioms
characterize the index uniquely, but an argument deriving Multiplication from the
other axioms is not known. Therefore we carry Multiplication along as an additional property that is extended from one level of generality to the next along with
everything else.
176
13.1
177
The axiom system for the fixed point index is introduced in two stages. This
section presents the first group of axioms, which describe the properties of the
index that concern a single space. Fix a metric space X. For a compact C X
let int C = C \ C be the interior of C, and let C = C \ int C be its topological
boundary.
Definition 13.1.1. An index admissible correspondence for X is an upper
semicontinuous correspondence F : C X, where C X is compact, that has no
fixed points in C.
There will be various indices, according to which sorts of correspondences are
considered. The next definition expresses the common properties of their domains.
Definition 13.1.2. An index base for X is a set of index admissible correspondences F : C X such that:
(a) f I whenever C X is compact and f : C X is an index admissible
continuous function;
(b) F |D I whenever F : C X is an element of I, D C is compact, and
F |D is index admissible.
For each integer m 0 an index base for Rm is given by letting I m be the set of
index admissible continuous functions f : C Rm .
We can now state the first batch of axioms.
Definition 13.1.3. Let I be an index base for X. An index for I is a function
X : I Z satisfying:
(I1) (Normalization1) If c : C X is a constant function whose value is an
element of int C, then
X (c) = 1.
(I2) (Additivity) If F : C X is an element of I, C1 , . . . , Cr are pairwise disjoint
compact subsets of C, and F P(F ) int C1 . . . int Cr , then
X (F ) =
X
i
X (F |Ci ).
178
13.2
Multiple Spaces
We now introduced two properties of the index that involve comparison across
different spaces. When we define an abstract notion of an index satisfying these
conditions, we need to require that the set of spaces is closed under the operations
that are involved in these conditions, so we require that the sets of spaces and
correspondences are closed under cartesian products.
Definition 13.2.1. An index scope S consists of a class of metric spaces SS and
an index base IS (X) for each X SS such that
(a) SS contains X X whenever X, X SS ;
(b) F F IS (X X ) whenever X, X SS , F IS (X), and F IS (X ).
179
Our first index scope S 0 has the collection of spaces SS 0 = {R0 , R1 , R2 , . . .} with
IS 0 (Rm ) = I m for each m. Of course (b) is satisfied by identifying Rm Rn with
Rm+n .
To understand the motivation for the following definition, first suppose that
SS , and that g : X X
and g : X
X are continuous. In this
X, X
circumstance it will be the case that g g and g g have the same index. We
and
would like to develop this idea in greater generality, for functions g : C X
g : C X, but for our purposes it is too restrictive to require that g(C) C and
C. In this way we arrive at the following definition.
g(C)
Definition 13.2.2. A commutativity configuration is a tuple
D,
E,
g)
(X, D, E, g, X,
are metric spaces and:
where X and X
X,
and D, D,
E, and E are compact;
(a) E D X, E D
and g C(D,
X) with g(E) int D
and g(E)
int D;
(b) g C(D, X)
(c) g g|E and g g|E are index admissible;
(d) g(F P(
g g|E )) = F P(g g|E ).
Before going forward, we should think through the details of what (d) means.
If x is a fixed point of g g|E , then g(x) is certainly a fixed point of g g, so it is a
Thus the inclusion
fixed point of g g|E if and only if g(x) E.
g(F P(
g g|E )) F P(g g|E )
holds if and only if
g(F P(
g g|E )) E.
()
()
Thus (d) holds if and only if both () and () hold, and by symmetry this is the
g g|E )).
case if and only if g(F P(g g|E )) = F P(
Definition 13.2.3. An index for an index scope S is a specification of an index
X for each X SS such that:
D,
E,
g) is a commutativity configuration
(I4) (Commutativity) If (X, D, E, g, X,
then
g g| ) IS (X),
with X, X SS , (E, g g|E ) IS (X), and (E,
E
X (
g g|E ) = X (g g|E ).
180
13.3
V1 V2 = im L,
V1 V3 = ker K,
and similarly for W . With suitably chosen bases the matrices of K and L have the
forms
0 L12 0 L14
0 K12 0 K14
0 K22 0 K24
and 0 L22 0 L24
0 0 0 0
0 0 0 0
0 0 0 0
0 0 0 0
181
g(F P(
g g|E )) int E
Since (g, g) 7 g g|E and (g, g) 7 g g|E are continuous, Theorem 5.2.1 and Lemma
4.5.10 imply that the set of such pairs is open.
182
Proof of Theorem 13.3.1. Uniqueness and (I1)-(I3) follow from Proposition 13.1.5,
so, we only need to prove that (I4) and (M) are satisfied.
E,
g) is a commutativity configuration. Lemma
Suppose that (Rm , D, E, g, Rm , D,
13.3.4 states that it remains a commutativity configuration if g and g are replaced
by functions in any sufficiently small neighborhood, and Lemma 13.3.3 implies that
g g|E and g g|E are continuous functions of (g, g), so, since we already know that
(I3) holds, it suffices to prove the equation of (I4) after such a replacement. Since
Rm ) (Proposition 10.2.7) we
the smooth functions are dense in C(D, Rm ) and C(D,
may assume that g and g are smooth. In addition, Sards theorem implies that
the regular values of IdE g g|E are dense, so after perturbing g by adding an
arbitrarily small constant, we can make it the case that 0 is a regular value. In the
same way we can add a small constant to g to make 0 a regular value of IdE g g|E ,
and if the constant is small enough it will still be the case that 0 is a regular value
of IdE g g|E .
Let x1 , . . . , xr be the fixed points of g g|E , and for each i let xi = g(xi ). Then
x1 , . . . , x
r are the fixed points of g
g |E . Let D1 , . . . , Dr be pairwise disjoint compact
1, . . . , D
r be pairwise disjoint open subsets
subsets of E with xi int Di , and let D
i ) with
of E with xi int Di . For each i let Ei be a compact subset of g 1(int D
1
xi int Ei , and let Ei be a compact subset of g (int Di ) with xi int Ei . It is easy
i, E
i , gi ) is a commutativity configuration.
to check that each (Rm , Di , Ei , gi , Rm , D
Recalling the relationship between the index and the degree, we have
g g|Ei ) = Rm (g g|Ei )
Rm (
because Proposition 13.3.2 gives
|I D(
g g)(xi )| = |I D
g (
xi )Dg(xi )| = |I Dg(xi )D
g (
xi )| = |I D(g g)(
xi )|.
Applying Additivity to sum over i gives the equality asserted by (I4).
13.4
Extension by Commutativity
The extension of the fixed point index to absolute neighborhood retracts was
first achieved by Felix Browder in his Ph.D. thesis Browder (1948), using the extension method described in this section. This extension method is, perhaps, the most
important application of Commutativity, but Commutativity is also sometimes useful in applications of the index, which should not be particularly surprising since
the underlying fact of linear algebra it embodies (Proposition 13.3.2) is already
nontrivial.
183
We
Throughout this section we will work with two fixed index scopes S and S.
say that S subsumes S if, for every X SS, we have X SS and IS(X) IS (X).
If this is the case, and is an index for S, then the restriction (in the obvious sense)
(It is easy to check that this is an automatic consequence
of to S is an index for S.
is an index for S,
then an extension to S is an
of the definition of an index.) If
E f (E) int D,
the supremum of the set of > 0 such that d(x , f (x )) > 2 whenever x
C \ (int E), x C, and d(x, x ) < ,
where d is the given metric for X. (Of course X has many metrics that give the
same topology. In contexts such as this we will implicitly assume that one has been
selected.)
Since f is continuous and admissible, narrowings of focus for f exist: continuity
implies the existence of an open neighborhood V of F P(f ) satisfying V f (V )
int C. Repeating this observation gives an open neighborhood W of F P(f ) satisfying W f (W ) V , and we can let D = V and E = W .
)-domination of C is
Let C be a compact subset of a metric space X. An (S,
: C C
()
IS ANR (X X)
184
Theorem 13.4.2. There is a unique index ANR for S ANR that extends 0 , and
ANR is multiplicative.
Proof. Theorem 7.6.4 implies that S 0 dominates S ANR , and S ANR evidently subsumes S 0 .
The rest of this section is devoted to the proof of Theorem 13.4.1. Before
proceeding, the reader should be warned that this is, perhaps, the most difficult
argument in this book. Certainly it is the most cumbersome, from the point of
view of the burden of notation, because the set up used to extend the index is
complex, and then several verifications are required in that setting. To make the
expressions somewhat more compact, from this point forward we will frequently
drop the symbol for composition, for instance writing rather than .
Lemma 13.4.3. Suppose X SS, f : C X is in IS(X), (D, E) is a narrowing
C,
, ) is an (S,
)-domination of U. Let
of focus for f , 0 < < (D,E) , and (X,
= 1 (D) and E = 1 (E). Then
D
D,
E,
| )
(X, D, E, f |D , X,
D
is a commutativity configuration.
Proof. We need to verify (a)-(d) of Definition 13.2.2. We have E D X with
X.
In addition D
and E are closed because is
D and E compact, so E D
Thus (a) holds.
continuous, so they are compact because they are subsets of C.
Of course f |D and |D are continuous. We have
((f (E))) U (f (E)) int D,
In addition, (E)
E int D. Thus (b) holds.
so (f (E)) 1 (int D) int D.
If x D \ int (E), then d(x, f (x)) > 2(D,E) > 2 and d(f (x), ((f (x)))) < ,
is a fixed
so x cannot be a fixed point of f . Thus F P(f |E ) int E. If x D
1
x) is a fixed point of f |D , so F P(f |D ) (int E)
point of f |D , then (
Thus (c) holds.
int E.
We now establish () and (). We have
((f (F P(f |D )))) = F P(f |D ) int E,
so
185
E1 = 11 (E),
2 = 1 (D),
D
2
E2 = 21 (E).
(2 f 2 | ).
(1 f 1 | ) =
E2
X2
E1
X1
X1
E1
X1
E1
X2
E2
X2
E2
Specifically, in view of (a) and (b) the first and third equalities follows from the
homotopy principle, while (c) permits an application of Commutativity that gives
the second equality.
For each t the composition 1 jt f 1 |E1 is well defined because
1 (E1 ) E
186
so
1 ) 21 (D) int D
2
2 f 1 (E
2 )) 11 (int D) = int D
1.
and 1 (2 (E
which is to say that () and () hold, which implies (d), completing the proof.
The hypotheses of the next result are mostly somewhat more general, but we
now need to assume that S dominates S.
Lemma 13.4.5. Assume that S dominates S. Let X be an element of SS , and let
f : C X be an element of IS (X). Suppose (D1 , E1 ) and (D2 , E2 ) are narrowings
1 , C1 , 1 , 1 ) and
of focus for f , 0 < 1 < (D1 ,E1 ) , 0 < 2 < (D2 ,E2 ) , and (X
2 , C2 , 2 , 2 ) are an (S,
1)-domination and an (S,
2 )-domination of C. Set
(X
1 = 1 (D1 ),
D
1
Then
E1 = 11 (E1 ),
2 = 1 (D2 ),
D
2
E2 = 21 (E2 ).
(2 f 2 | ).
(1 f 1 | ) =
E2
X2
E1
X1
for arbitrarily small 2 , and if we apply the lemma above to this domination and
the given one we find that it suffices to prove the result with the given domination
replaced by this one. This means that we may assume that 2 is as small as need
be, and in particular we may assume that 2 < (D1 ,E1 ) . Now Additivity implies
that
(2 f 2 | ) =
(2 f 2 | 1
2 (E1 ) ),
X2
E2
X2
which means that it suffices to establish the result with D2 and E2 replaced by D1
and E1 , which is the case established in the lemma above.
Proof of Theorem 13.4.1. Since S dominates S, the objects used to define exist,
and the last result implies that the definition of does not depend on the choice
C,
, ). We now verify that satisfies (I1)-(I4) and (M).
of (D, E), , and (X,
187
Normalization:
gives
If f is a constant function, then so is f , so Normalization for
(f ) = 1.
X (f ) =
X
Additivity:
Suppose that F P(f ) int C1 . . . int Cr where C1 , . . . , Cr C are compact
and pairwise disjoint. For each j = 1, . . . , r choose open sets Dj D Cj and
Ej E Cj such that (Dj , Ej ) is a narrowing of focus for (Cj , f |Cj ). In view of
Lemma 13.4.5 we may assume that < (Dj ,Ej ) for all j. It is easy to see that for
C,
|C , ) is an (S,
)-domination of Cj . For each j let E = 1 (Ej ).
each j, (X,
j
j
gives
Additivity for
X
X
(f |E ) =
(f | ) =
X (f |Cj ).
X (f ) =
X
E
X
j
j
Continuity:
It is easy to see that if f : C X that are sufficiently close to f , then (D, E) is
C,
, ) is a (S,
)-domination of C. Since
a narrowing of focus for (C, f ), and (X,
gives
f 7 f is continuous (Propositions 5.5.2 and 5.5.3) Continuity for
(f ) = X (f )
(f ) =
X (f ) =
X
X
when f is sufficiently close to f .
Commutativity:
Suppose that (X1 , C1 , D1 , g1 , X2 , C2 , D2 , g2 ) is a commutativity configuration
with X1 , X2 SS . Replacing D1 and D2 with smaller open neighborhoods of
F P(g2 g1 ) and F P(g1 g2 ) if need be, we may assume that
D1 g2 g1 (D1 ) C1
and D2 g1 g2 (D2 ) C2 .
E1 = 11 (E1 ),
2 = 21 (D2 ),
D
E2 = 21 (E2 ).
188
Here the first and fifth equality are from the definition of , the second and fourth
and the third is from Commutativity for .
In
are implied by Continuity for ,
order for this to work it must be the case that all the compositions in this calculation
are well defined, in the sense that the image of the first function is contained in the
domain of the second function, the homotopies
t 7 1 g2 jt g1 1 |E1
and t 7 2 g1 ht g2 2 |E2
1 X
2 , C1 C2 , 1 2 , 1 2 ) is a (S,
)-domination
Let = max{1 , 2 }. Then (X
of C1 C2 . It is also easy to check that (D1 D2 ,E1 E2 ) max{(D1 ,E1 ) , (D2 ,E2 ) }, so
< (D1 D2 ,E1 E2 ) . Therefore Lemma 13.4.3 implies that the validity of the first
equality in
((1 f1 1 2 f2 2 )| )
X1 X2 (f1 f2 ) =
X1 X2
E 1 E 2
(2 f2 2 | ) = X1 (f1 ) X2 (f2 )
(1 f1 1 | )
=
E2
X2
E1
X1
and the third is the definition
the second one is an application of Multiplication for
of .
then is the unique extension of
to
We now prove that if S subsumes S,
S. Consider X SS and (C, f ) IS(X). For any > 0, (X, C, IdC , IdC ) is
)-domination of C. For any narrowing of focus (D, E) equation () gives
an (S,
X (f |E ) and Additivity for
gives
X (f |E ) =
X (f ). Thus extends
X (f ) =
.
necessarily agree everywhere because, by
Two indices for S that restrict to
Continuity and Commutativity, () holds in the circumstances described in the
statement of Theorem 13.4.1.
13.5
189
Extension by Continuity
This section extends the index from continuous functions to upper semicontinuous contractible valued correspondences. As in the last section, we describe the
extension process abstractly, thereby emphasizing the aspects of the situation that
drive the argument.
Definition 13.5.1. If I and I are index bases for a compact metric space X, we
say that I approximates I if:
(E1) If C, D X are open with D C, then I C(D, X) is dense in
{ F |D : F I U(C, X) and F |D is index admissable }.
(E2) If C, D X are open with D C, F I U(C, X), and A C X is a
neighborhood of Gr(F ), then there is a neighborhood B D X of Gr(F |D )
such that any two functions f, f C(D, X) with Gr(f ), Gr(f ) B are the
endpoints of a homotopy h : [0, 1] C(D, X) with Gr(ht ) A for all t.
It would be simpler if, in (E1) and (E2), we could have V = U, but unfortunately
Theorem 9.1.1 is not strong enough to justify working with such a definition.
Definition 13.5.2. If S and S are two index scopes with SS = SS , then S approximates S) if, for each X SS, IS(X) approximates IS (X), and
190
Proposition 13.5.4. Suppose I and I are index bases for a compact metric space
X : I Z there is a unique
X, and I approximates I. Then for any index
index X : I Z such that for each open C X with compact closure, each F
I U(C, X), and each open D with F P(F ) D and D C, there is a neighborhood
X (f ) for all f E C(D, X) I.
h0 = f , h1 = f , and
}
Gr(ht ) (C X) \ { (x, x) : x C \ D
for all t. Since restriction to a compact subdomain is a continuous
(Lemma 5.3.1) we may replace B with a smaller neighborhood of F |D
Gr(f |D ) B whenever Gr(f ) B. For such an f Additivity gives
X (f | ) as desired.
D
It remains to show that (I1)-(I3) are satisfied.
operation
to obtain
X (f ) =
Normalization:
If c is a constant function, we can take c itself as the approximation used to
X.
define X (c), so Normalization for X follows from Normalization for
Additivity:
Consider F I with domain C. Let C1 , . . . , Cr be disjoint open subsets of C
whose union contains F P(F ). Let D1 , . . . , Dr be open subsets of C with D1
191
X (F |Ci ).
i
i
192
Then the definition of an index scope implies that F F IS (XX ). Choose open
sets D and D with F P(F ) D, D C, F P(F ) D , and D C . As above, we
can find neighborhoods B U(D, X), B U(D , X ), and D U(DD , X X ),
X (f ) for all
of F |D , F |D , and (F F )|DD respectively, such that X (F ) =
S
of correspondences is a continuous operation (this is Lemma 5.3.4) we may replace
B and B with smaller neighborhoods to obtain F F D for all F B and
F B . Assumption (E1) implies that there are
f B IS(X) C(D, X) and f B IS(X ) C(D , X ).
The definition of an index scope implies that f f IS(X X ), and Multiplication
now gives
(I4) for
X (f )
X (f ) = X (F ) X (F ).
XX (f f ) =
XX (F F ) =
Part III
Applications and Extensions
193
Chapter 14
Topological Consequences
This chapter is a relaxing and refreshing change of pace. Instead of working very
hard to slowly build up a toolbox of techniques and specific facts, we are going to
harvest the fruits of our earlier efforts, using the axiomatic description of the fixed
point index, and other major results, to quickly derive a number of quite famous
results. In Section 14.1 we define the Euler characteristic, relate it to the Lefschetz
fixed point theorem, and then describe the Eilenberg-Montgomery as a special case.
For two general compact manifolds, the degree of a map from one to the other
is a rather crude invariant, in comparison with many others that topologists have
defined. Nevertheless, when the range is the m-dimensional sphere, the degree
is already a complete invariant in the sense that it classifies functions up to
homotopy: if M is a compact m-dimensional manifold that is connected, and f
and f are functions from M to the m-sphere of the same degree, then f and f
are homotopic. This famous theorem, due to Hopf, is the subject of Section 14.2.
Section 12.4 proves a simple result asserting that the degree of a composition of two
functions is the products of their degrees.
Section 14.3 presents several other results concerning fixed points and antipodal
maps of a map from a sphere to itself. Some of these are immediate consequences
of index theory and the Hopf theorem, but the Borsuk-Ulam theorem requires a
substantial proof, so it should be thought of as a significant independent fact of
topology. It has many consequences, including the fact that spheres of different
dimensions are not homeomorphic.
In Section 14.4 we state and prove the theorem known as invariance of domain.
It asserts that if U Rm is open, and f : U Rm is continuous and injective, then
the image of f is open, and the inverse is continuous. One may think of this as a
purely topological version of the inverse function theorem, but from the technical
point of view it is much deeper.
If a connected set of fixed points has a nonzero index, it is essential. This raises
the question of whether a connected set of fixed points of index zero is necessarily
inessential. Section 14.5 presents two results of this sort.
194
14.1
195
The definition of the Euler characteristic, and Eulers use of it in the analyses
of various problems, is often described as the historical starting point of topology
as a branch of mathematics. In popular expositions the Euler characteristic of a 2dimensional manifold M is usually defined by the formula (M) := V E +F where
V , E, and F are the numbers of vertices, edges, and 2-simplices in a triangulation
of M. Our definition is:
Definition 14.1.1. The Euler characteristic (X) of a compact ANR X is
X (IdX ).
Here is a sketch of a proof that our definition of (M) agrees with Eulers
when M is a triangulated compact 2-manifold. We deform the identity function
slightly, achieving a function f : M M defined as follows. Each vertex of the
triangulation is mapped to itself by f . Each barycenter of an edge is mapped to
itself, and the points on the edge between the barycenter and either of the vertices
of the edge are moved toward the barycenter. Each barycenter of a two dimensional
simplex is mapped to itself. If x is a point on the boundary of the 2-simplex, the
line segment between x and the barycenter is mapped to the line segment between
f (x) and the barycenter, with points on the interior of the line segment pushed
toward the barycenter, relative to the affine mapping. It is easy to see that the only
fixed points of f are the vertices and the barycenters of the edges and 2-simplices.
Eulers formula follows once we show that the index of a vertex is +1, the index of
the barycenter of an edge is 1, and the index of the barycenter of a 2-simplex is +1.
We will not give a detailed argument to this effect; very roughly it corresponds to
the intuition that f is expansive at each vertex, compressive at the barycenter
of each 2-simplex, and expansive in one direction and compressive in another at the
barycenter of an edge.
Although Euler could not have expressed the idea in modern language, he certainly understood that the Euler characteristic is important because it is a topological invariant.
Theorem 14.1.2. If X and X are homeomorphic compact ANRs, then
(X) = (X ).
Proof. For any homeomorphism h : X X , Commutativity implies that
(X) = X (IdX ) = X (IdX h1 h) = X (h IdX h1 ) = X (IdX ) = (X ).
196
whether a triangulation exists. That a two dimensional compact manifold is triangulable was not proved until the 1920s, by Rado. In the 1950s Bing and Moise
proved that compact three dimensional manifolds are triangulable, and a stream
of research during this same general period showed that smooth manifolds are triangulable, but in general a compact manifold need not have a triangulation. For
simplicial complexes topological invariance would follow from invariance under subdivision, which can be proved combinatorically, and the Hauptvermutung, which
was the conjecture that any two simplicial complexes that are homeomorphic have
subdivisions that are combinatorically isomorphic. This conjecture was formulated
by Steinitz and Tietze in 1908, but in 1961 Milnor presented a counterexample, and
in the late 1960s it was shown to be false even for triangulable manifolds.
The Lefschetz fixed point theorem is a generalization Brouwers theorem
that was developed by Lefschetz for compact manifolds in Lefschetz (1923, 1926)
and extended by him to manifolds with boundary in Lefschetz (1927). Using quite
different methods, Hopf extended the result to simplicial complexes in Hopf (1928).
Definition 14.1.3. If X is a compact ANR and F : X X is an upper semicontinuous contractible valued correspondence, the Lefschetz number of F is X (F ).
Theorem 14.1.4. If X is a compact ANR, F : X X is an upper semicontinuous
contractible valued correspondence and X (F ) 6= 0, then F P(F ) 6= .
Proof. When F P(F ) = two applications of Additivity give
(F | ) = (F ) = (F | ) + (F |).
In Lefschetz originally formulation the Lefschetz number of a function was defined using algebraic topology. Thus one may view the Lefschetz fixed point theorem
as a combination of the result above and a formula expressing the Lefschetz number
in terms of homology.
In the Kakutani fixed point theorem, the hypothesis that the correspondence is
convex valued cries out for generalization, because convexity is not a topological concept that is preserved by homeomorphisms of the space. The Eilenberg-Montgomery
theorem asserts that if X is a compact acyclic ANR, and F : X X is an upper
semicontinuous acyclic valued correspondence, then F has a fixed point. Unfortunately it would take many pages to define acyclicity, so we will simply say that
acyclicity is a property that is invariant under homeomorphism, and is weaker than
contractibility. The known examples of spaces that are acyclic but not contractible
are not objects one would expect to encounter in nature, so it seems farfetched
that the additional strength of the Eilenberg-Montgomery theorem, beyond that of
the result below, will ever figure in economic analysis.
Theorem 14.1.5. If X is a nonempty compact absolute retract and F : X X
is an upper semicontinuous contractible valued correspondence, then F has a fixed
point.
197
Proof. Recall (Proposition 7.5.3) that an absolute retract is an ANR that is contractible. Theorem 9.1.1 implies that F can be approximated in the sense of
Continuity by a continuous function, so X (F ) = X (f ) for some continuous
f : X X. Let c : X [0, 1] X be a contraction. Then (x, t) 7 c(f (x), t) (or
(x, t) f (c(x, t))) is a homotopy between f and a constant function, so Homotopy
[fix this] and Normalization imply that X (f ) = 1. Now the claim follows from the
last result.
14.2
Two functions that are homotopic may differ in their quantitative features, but
from the perspective of topology these differences are uninteresting. Two functions
that are not homotopic differ in some qualitative way that one may hope to characterize in terms of discrete objects. A homotopy!invariant may be thought of
as a function whose domain is the set of homotopy classes; equivalently, it may be
thought of as a mapping from a space of functions that is constant on each homotopy class. A fundamental method of topology is to define and study homotopy
invariants.
The degree is an example: for compact manifolds M and N of the same dimension it assigns an integer to each continuous f : M N, and if f and f are
homotopic, then they have the same degree. There are a great many other homotopy invariants, whose systematic study is far beyond our scope. In the study of
such invariants, one is naturally interested in settings in which some invariant (or
collection of invariants) gives a complete classification, in the sense that if two functions are not homotopic, then the invariant assigns different values to them. The
prototypical result of this sort, due to Hopf, asserts that the degree is a complete
invariant when N is the m-sphere.
Theorem 14.2.1 (Hopf). If M is an m-dimensional compact connected smooth
manifold, then two maps f, f : M S m are homotopic if and only if deg(f ) =
deg(f ).
We provide a rather informal sketch of the proof. Since the ideas in the argument
are geometric, and easily visualized, this should be completely convincing, and little
would be gained by adding more formal details of particular constructions.
We already know that two homotopic functions have the same degree, so our
goal is to show that two functions of the same degree are homotopic. Consider a
particular f : M S m . The results of Section 10.7 imply that CS (M, S m ) is locally
path connected, and that C (M, S m ) is dense in this space, so f is homotopic to a
smooth function. Suppose that f is smooth, and that q is a regular value of f . (The
existence of such a q follows from Sards theorem.) The inverse function theorem
implies that if D is a sufficiently small disk in S m centered at q, then f 1 (D) is a
collection of pairwise disjoint disks, each containing one element of f 1 (q).
Let q be the antipode of q in S m . (This is q when S m is the unit sphere
centered at the origin in Rm+1 .) Let j : S m [0, 1] S m be a homotopy with
j0 = IdS m that stretches D until it covers S m , so that j1 maps the boundary of D
and everything outside D to q . Then f = j0 f is homotopic to j1 f .
198
kx e1 k < 1,
x e1 ,
f (x) = x (e1 ) 2(hx, e1 i he1 , e1 i)e1 , kx (e1 )k < 1,
q ,
otherwise.
q ,
otherwise.
Of course the first two expressions agree when x1 = 0, so this is well defined and
continuous, and h1 (x) = q for all x.
In preparation for an application of the Hopf theorem, we introduce an important
concept from topology. If X is a topological space and A X, the pair (X, A)
has the homotopy!extension property if, for any topological space Y and any
function g : (X {0}) (A [0, 1]) Y , there is a homotopy h : X [0, 1] Y
such that is an extension of g: h(x, 0) = g(x, 0) for all x X and h(x, t) = g(x, t)
for all (x, t) A [0, 1].
199
Lemma 14.2.2. The pair (X, A) has the homotopy extension property if and only
if (X {0}) (A [0, 1]) is a retract of X [0, 1].
Proof. If (X, A) has the homotopy extension property, then the inclusion map from
(X {0}) (A [0, 1]) to X [0, 1] has a continuous extension to all of X [0, 1],
which is to say that there is a retraction. On the other hand, if r is a retraction,
then for any g there is continuous extension h = g r.
We will only be concerned with the example given by the next result, but it
is worth noting that this concept takes on greater power when one realizes that
(X, A) has the homotopy extension property whenever X is a simplicial complex
and A is a subcomplex. It is easy to prove this if there is only one simplex in X
that is not in A; either the boundary of is contained in A, in which case there
is an argument like the proof of the following, or it isnt, and another very simple
construction works. The general case follows from induction because if (X, A) and
(A, B) have the homotopy extension property, then so does (X, B). To show this
suppose that g : (X {0}) (B [0, 1]) Y is given. There is a continuous
extension h : A [0, 1] Y of the restriction of g to (A {0}) (B [0, 1]). The
extension of h to all of (X {0}) (A [0, 1]) defined by setting h|X{0} = g|X{0}
is continuous because it is continuous on X {0} and A [0, 1], both of which
are closed subsets of X [0, 1] (here the requirement that A is closed finally shows
up) and since (X, A) has the homotopy extension property this h can be further
extended to all of X [0, 1].
Lemma 14.2.3. The pair (D m , S m1 ) has the homotopy extension property.
Proof. There is an obvious retraction
r : D m [0, 1] (D m {0}) (S m1 [0, 1])
defined by projecting radially from (0, 2) Rm R.
We now relate the degree of a map from D m to Rm with what may be thought
of as the winding number of the restriction of the map to S m1 .
Theorem 14.2.4. If f : D m Rm is continuous, 0
/ f (S m1 ), and f : S m1
m1
S
is the function x 7 f (x)/kf (x)k, then deg0 (f ) = deg(f).
Proof. For k Z let fk : D m Rm be the map
(r cos , r sin , x3 , . . . , xm ) 7 (r cos k, r sin k, x3 , . . . , xm ).
It is easy to see that deg0 (fk ) = k = deg(f |S m1 ).
Now let k = deg(f). The Hopf theorem implies that there is a homotopy
: S m1 [0, 1] S m1 with h
0 = f and h
1 = fk |S m1 . Let h : S m1 [0, 1] Rm
h
be the homotopy with h0 = f |S m1 and h1 = fk |S m1 given by
t),
h(x, t) = (1 t)kf (x)k + t h(x,
200
14.3
Insofar as spheres are the simplest nontrivial (where, in effect, this means
noncontractible) topological spaces, it is entirely natural that mathematicians would
quickly investigate the application of degree and index theory to these spaces, and
to maps between them. There are many results coming out of this research, some
of which are quite famous.
Our discussion combines some purely topological reasoning with analysis based
on concrete examples, and for the latter it is best to agree that
S m := { x Rm+1 : kxk = 1 }.
Some of our arguments involve induction on m, and for this purpose we will regard
S m1 as a subset of S m by setting
S m1 = { x S m : xm+1 = 0 }.
Let am : S m S m be the function
am (x) = x.
Two points x, y S m are said to be antipodal if y = am (x). Regarded topologically, am is a fixed point free local diffeomorphism whose composition with itself
is IdS m , and one should expect that all the topological results below involving am
and antipodal points should depend only on these properties, but we will not try to
demonstrate this (the subject is huge, and our coverage is cursory) instead treating
am as an entirely concrete object.
Let
Em = { (x, y) S m S m : y 6= am (x) }.
There is a continuous function rm : Em [0, 1] S m given by
rm (x, y, t) :=
tx + (1 t)y
.
ktx + (1 t)yk
201
deg(f1,d ) = d.
Now observe that f1,1 is homotopic to a map without fixed points, while for
d 6= 1 the fixed points of f1,d are the points
2k
2k
, sin d1
(k = 0, . . . , d 2).
cos d1
If d > 1, then motion in the domain is translated by f1,d into more rapid motion
in the range, so the index of each fixed point is 1. When d < 1, f1,d translates
motion in the domain into motion in the opposite direction in the range, so the
index of each fixed point is 1. Combining these facts, we conclude that
(f1,d ) = 1 d,
which establishes the result when m = 1.
Let em+1 = (0, . . . , 0, 1) Rm+1 . Then
S m = { x + em+1 : x S m1 , 0, 2 + 2 = 1 }.
202
203
204
Sards theorem implies that some q S m1 is a regular value of both g and g|M .
Theorem 12.2.1 implies that degq (g|M ) = 0, so (g|M )1 (q ) has an even number of
elements. Evidently g maps the boundary of each Cp diffeomorphically onto S m1 ,
so each such boundary contains exactly one element of (g|M )1 (q ). In addition,
j maps antipodal points of S m \ {q, q} to antipodal ponts of S m1 , so g|S m1 is
antipodal, and our induction hypothesis implies that (g|M )1 (q ) S m1 has an
odd number of elements. Therefore the number of components of f 1 (D D )
contained in S+m is odd, as desired.
The hypotheses can be weakened:
Corollary 14.3.7. If the map f : S m S m satisfies f (p) 6= f (p) for all p, then
the degree of f is odd.
Proof. This will follow from the last result once we have shown that f is homotopic to an antipodal map. Let h : S m [0, 1] S m be the homotopy h(p, t) =
rm (f (p), f (p), 2t). The hypothesis implies that this is well defined, and h1 is
antipodal.
This result has a wealth of geometric consequences.
Theorem 14.3.8 (Borsuk-Ulam Theorem). The following are true:
(a) If f : S m Rm is continuous, then there is a p S m such that f (p) =
f (am (p)).
(b) If f : S m Rm is continuous and antipodal, then there is a p S m such that
f (p) = 0.
(c) There is no continuous antipodal f : S m S m1 .
(d) There is no continuous g : D m = { (y1, . . . , ym , 0) Rm+1 : kyk 1 } S m1
such that g|S m1 is antipodal.
205
(e) Any cover F1 , . . . , Fm+1 of S m by m + 1 closed sets has a least one set that
contains a pair of antipodal points.
(f ) Any cover U1 , . . . , Um+1 of S m by m + 1 open sets has a least one set that
contains a pair of antipodal points.
Proof. We think of Rm as S m with a point removed, so a continuous f : S m Rm
amounts to a function from S m to itself that is not surjective, and whose degree is
consequently zero. Now (a) follows from the last result.
Suppose that f : S m Rm is continuous and f (p) = f (p). If f is also
antipodal, then f (p) = f (p) so f (p) = 0. Thus (a) implies (b).
Obviously (b) implies (c).
Let : p 7 (p1 , . . . , pm , 0) be the standard projection from Rm+1 to Rm . As
in the proof of Theorem 14.3.6 let S+m and Sm be the Northern and Southern
hemispheres of S m . If g : D m S m1 was continuous and antipodal, we could
define a continuous and antipodal f : S m S m1 by setting
(
g((p)),
p S+m ,
f (p) =
g((am (p))), p Sm .
Thus (c) implies (d).
Suppose that F1 , . . . , Fm+1 is a cover of S m by closed sets. Define f : S m Rm
by setting
f (p) = d(x, F1 ), . . . , d(x, Fm )
where d(x, x ) = kx x k is the usual metric for Rm+1 . Suppose that f (p) =
f (p) = y. If yi = 0, then p, p Fi , and if all the components of y are nonzero,
then p, p Fm+1 . Thus (a) implies (e).
Suppose U1 , . . . , Um+1 is a cover of S m by open sets and > 0. For i = 1, . . . , m+
1 set Fi := { p S m : d(p, S m \ Ui ) }. Then each Fi is a closed subset of Ui ,
and these sets cover S m if is sufficiently small. Thus (e) implies (f).
In the argument above we showed that (a) (b) (c) (d) and (a) (e)
(f). There are also easy arguments for the implications (d) (c) (b) (a) and
(f) (e) (c), so (a)-(f) are equivalent in the sense of each being an elementary
consequence of each other. The proofs that (d) (c) and (c) (b) are obvious
and can be safely left to the reader. To show that (b) (a), for a given continuous
f : S m Rm we apply (b) to f f am . To show that (f) (e) observe that if
F1 , . . . , Fm+1 are closed and cover S m , then for each n the sets U1/n (Fi ) are open
and cover S m , so there is a pn with pn , pn U1/n (Fi ) for some i. Any limit point
of the sequence {pn } has the desired property.
The proof that (e) (c) is more interesting. Consider an m-simplex that is
embedded in D m with the origin in its interior. Let F1 , . . . , Fm+1 be the radial
projections of the facets of the simplex onto S m1 . These sets are closed and cover
S m1 , and since each facet is separated from the origin by a hyperplane, each Fi
does not contain an antipodal pair of points. If f : S m S m1 is continuous,
then f 1 (F1 ), . . . , f 1 (Fm+1 ) are a cover of S m by closed sets, and (e) implies the
206
existence of p, p f 1 (Fi ) for some i. If f was also antipodal, then f (p), f (p) =
f (p) Fi , which is impossible.
As a consequence of the Borsuk-Ulam theorem, the following obvious fact is
actually highly nontrivial.
Theorem 14.3.9. Spheres of different dimensions are not homeomorphic.
Proof. If k < m then, since S k can be embedded in Rm , part (a) of the Borsuk-Ulam
theorem implies that a continuous function from S m to S k cannot be injective.
14.4
Invariance of Domain
The main result of this section, invariance of domain, is a famous result with
numerous applications. It can be thought of as a purely topological version of the
inverse function theorem. However, before that we give an important consequences
of the Borsuk-Ulam theorem for Euclidean spaces.
Theorem 14.4.1. Euclidean spaces of different dimensions are not homeomorphic.
Proof. If k 6= m and f : Rk Rm was a homeomorphism, for any sequence {xj } in
Rk with {xj } the sequence {f (xj )} could not have a convergent subsequence,
so kf (xj )k . Identifying Rk and Rm with S k \ {ptk } and S m \ {ptm }, the
extension of f to S k given by setting f (ptk ) = ptm would be continuous, with a
continuous inverse, contrary to the last result.
The next two lemmas develop the proof of this sections main result.
Lemma 14.4.2. Suppose S+m is the Northern hemisphere of S m , f : S+m S m
is a map such that f |S m1 is antipodal, and p S+m \ S m1 is a point such that
p
/ f (S+m ) and p
/ f (S m1 ). Then degp (f ) is odd.
Proof. Let f : S m S m be the extension of f given by setting f(p) = f (p)
when pm+1 < 0. Clearly f is continuous and antipodal, so its degree is odd. The
hypotheses imply that f1 (p) S+m \ S m1 , and that f is degree admissible over p,
so Additivity implies that degp (f ) = degp (f).
Lemma 14.4.3. If f : D m Rm is injective, then degf (0) (f ) is odd, and f (D m )
includes a neighborhood of f (0).
Proof. Replacing f with x 7 f (x) f (0), we may assume that f (0) = 0. Let
h : D m [0, 1] Rm be the homotopy
x
h(x, t) := f ( 1+t
) f ( tx
).
1+t
207
The next result is quite famous, being commonly regarded as one of the major
accomplishments of algebraic topology. As the elementary nature of the assertion
suggests, it is applied quite frequently.
Theorem 14.4.4 (Invariance of Domain). If U Rm is open and f : U Rm
is continuous and injective, then f (U) is open and f is a homeomorphism onto its
image.
Proof. The last result can be applied to a closed disk surrounding any point in the
domain, so for any open V U, f (V ) is open. Thus f 1 is continuous.
14.5
208
= : V M M
is a C r1 embedding, where : T M M is the projection. Let V =
(V ).
Let Y0 = { p C : (p, f (p)) V }; of course this is an open set containing
F P(f ). Let Y1 and Y2 be open sets such that F P(f ) Y2 , Y 2 Y1 , Y 1 Y0 ,
and Y2 is path connected. (Such a Y2 can be constructed by taking a finite union
of images of D m under C r parameterizations.) We can define a vector field on a
neighborhood of Y0 by setting
(p) =
1 (p, f (p)).
Proposition 11.5.2 and Corollary 10.2.5 combine to imply that there is a vector field
on Y0 with image contained in
1 (W ) that agrees with on Y0 \ Y1 , is C r1 on
Y2 , and has only regular equilibria, all of which are in Y2 . The number of equilibria
is necessarily finite, and we may assume that, among all the vector fields on Y0 that
agree with on Y0 \ Y1 , are C r1 on Y2 , and have only regular equilibria in Y2 ,
minimizes this number. If has no equilibria, then we may define a continuous
function f : C X without any fixed points whose graph is contained in W by
209
(This is obvious, but painful to prove formally, and in addition the case m = 1
requires special treatment. A formal verification would do little to improve the
readers understanding, so we omit the details.) Applying the tubular neighborhood
theorem, this path can be used to construct a C r parameterization : Z U where
Z Rm is a a neighborhood of D m .
Let g : Z Rm be defined by setting g(x) = D(x)1 (x) . Proposition 14.5.1
gives a continuous function g : Z Rm \ {0} that agrees with g on the closure of
Z \ D m . We extend g to all of Z by setting g(x) = g(x) if x
/ D m . Define a new
vector field on (Z) by setting
= D(1 (p))
(p)
g (1 (p)).
There are two final technical points. In order to insure that (p)
1 (W ) for
r
m
all p we can first multiplying g by a C function : D (0, 1] that is identically
1 on Z \ D m and close to zero in the interior of D m outside of some neighborhood of
S m1 . We can also use Proposition 11.5.2 and Corollary 10.2.5 to further perturb
to make is C r1 without introducing any additional equilibria. This completes the
construction, thereby arriving at a contradiction that completes the proof.
Economic applications call for a version of the result for correspondences. Ideally
one would like to encompass contractible valued correspondences in the setting of
a manifold, but the methods used here are not suitable. Instead we are restricted
to convex valued correspondences, and thus to settings where convexity is defined.
Theorem 14.5.3. If X Rm is compact and convex, C X is compact, F :
C X is an index admissible upper semicontinuous convex valued correspondence,
(F ) = 0, and F P(F ) is connected, then F is inessential.
Caution: The analogous result does not hold for essential sets of Nash equilibria,
which are defined by Jiang (1963) in terms of perturbations of the games payoffs.
Hauk and Hurkens (2002) give an example of a game with a component of the set
of Nash equilibria that has index zero but is robust with respect to perturbations
of payoffs.
Proof. Let W C X be an open set containing the graph of F . We will show that
there is a continuous f : C X with Gr(f ) W and F P(f ) = . Let x0 be a point
in the interior of X, let h : X [0, 1] X be the contraction h(x, t) = (1t)x+tx0 ,
and for t [0, 1] let ht F be the correspondence x 7 ht (F (x)). This correspondence
is obviously upper semicontinuous and convex valued, and Gr(ht F ) W for small
t > 0, so it suffices to prove the result with F replaced by ht F for such t. Therefore
we may assume that the image of F is contained in the interior of X.
For each x F P(F ) we choose convex neighborhoods Yx C of x and Zx X
of F (x) such that Yx Zx and Yx Zx W . Choose x1 , . . . , xk such that F P(F )
Yx1 . . . Yxk , and let
Y 0 = Y x1 . . . Y xk
Note that for all (x, y) Z0 , Z0 contains the line segment { (x, (1 t)y + tx) }. Let
Y1 and Y2 be open subsets of C with F P(F ) Y2 , Y 2 Y1 , Y 1 Y0 , and Y2 is
210
There are two final technical points. In order to insure that k(x)k
< for all
m
p we can first multiply g by a C function : D (0, 1] that is identically 1
on T \ D m and close to zero in the interior of D m outside of some neighborhood of
S m1 . We can also use Proposition 11.5.2 and Corollary 10.2.5 to further perturb
to make it C without introducing any additional zeros. We can now define
Chapter 15
Vector Fields and their Equilibria
Under mild technical conditions, explained in Sections 15.1 and 15.2, a vector
field on a manifold M determines a dynamical system. That is, there is a function
: W M, where W M R is a neighborhood of W {0}, such that the
derivative of at (p, t) W , with respect to time, is (p,t) . In this final chapter
we develop the relationship between the fixed point index and the stability of rest
points, and sets of rest points, of such a dynamical system.
In addition to the degree and the fixed point index, there is a third expression of
the underlying mathematical principle for vector fields. In Section 15.3 we present
an axiomatic description of the vector field index, paralleling our axiom systems
for the degree and fixed point index. Existence and uniqueness are established by
showing that the vector field index of |C , for suitable compact C M, agrees
with the fixed point index of (, t)|C for small negative t. Since we are primarily
interested in forward stability, it is more to the point to say that the fixed point
index of (, t)|C for small positive t agrees with the vector field index of |C .
The notion of stability we focus on, asymptotic stability, has a rather complicated definition, but the intuition is simple: a compact set A is asymptotically
stable if the trajectory of each point in some neighborhood of A is eventually drawn
into, and remains inside, arbitrarily small neighborhoods of A. In order to use the
fixed point index to study stability, we need to find some neighborhood of such an A
that is mapped into itself by (, t) for small positive t. The tool we use to achieve
this is the converse Lyapunov theorem, which asserts that if A is asymptotically
stable, then there is a Lyapunov function for that is defined on a neighborhood of
A. Unlike the better known Lyapunov theorem, which asserts that the existence of
a Lyapunov function implies asymptotic stability, the converse Lyapunov theorem
is a more recent and difficult result. We prove a version of it that is sufficient for
our needs in Section 15.5.
Once all this background material is in place, it will not take long to prove the
culminating result, that if A is a asymptotically stable, and an ANR, then the vector
field index of is the Euler characteristic of A. This was proved in the context of
a game theoretic model by Demichelis and Ritzberger (2003). The special case of A
being a singleton is a prominent result in the theory of dynamical systems, due to
Krasnoselski and Zabreiko (1984): if an isolated rest point is asymptotically stable
for , then the vector field index of that point for is 1.
211
212
Paul Samuelson advocated a correspondence principle in two papers Samuelson (1941, 1942) and his famous book Foundations of Economic Analysis Samuelson
(1947). The idea is that the stability of an economic equilibrium, with respect to
natural dynamics of adjustment to equilibrium, implies certain qualitative properties of the equilibriums comparative statics. There are 1-dimensional settings
in which this idea is regarded as natural and compelling, but Samuelsons writings discuss many examples without formulating it as a general theorem, and its
nature and status in higher dimensions has not been well understood; Echenique
(2008) provides a concise summary of the state of knowledge and related literature.
The book concludes with an explanation of how the Krasnoselski-Zabreiko theorem
allows the correspondence principle to be formulated in a precise and general way.
15.1
We begin with a review of the theory of ordinary differential equations in Euclidean space. Let U Rm be open, and let z : U Rm be a function, thought
of as a vector field. A trajectory of z is a C 1 function : (a, b) U such that
(s) = z(s) for all s. Without additional assumptions the dynamics associated
with z need not be deterministic: there can be more than one trajectory for the
vector field satisfying an initial condition that specifies the position of the trajectory
at a particular moment. For example, suppose that m = 1, U = R, and
(
0,
t 0,
z(t) =
2 t, t > 0.
Then for any s0 there is a trajectory s0 : R M given by
(
0,
s s0 ,
s0 (s) =
2
(s s0 ) , s > s0 .
For most purposes this sort of indeterminacy is unsatisfactory, so we need to find
a condition that implies that for any initial condition there is a unique trajectory.
Let (X, d) and (X , d ) be metric spaces. A function f : X X is Lipshitz if
there is a constant L > 0 such that
d (f (x), f (y)) Ld(x, y)
for all x, y X. We say that f is locally Lipschitz if each x X has a neighborhood U such that f |U is Lipschitz. The basic existence-uniqueness result for
ordinary differential equations is:
Theorem 15.1.1 (Picard-Lindelof Theorem). Suppose that U Rm is open, z :
U Rm is locally Lipschitz, and C U is compact. Then for sufficiently small
> 0 there is a unique function F : C (, ) U such that for each x C,
F (x, 0) = x and F (x, ) is a trajectory of z. In addition F is continuous, and if z
is C s (1 s ) then so is F .
213
Due to its fundamental character, a detailed proof would be out of place here,
but we will briefly describe the central ideas of two methods. First, for any > 0
one can define a piecewise linear approximate solution going forward in time by
setting F (x, 0) = x and inductively applying the equation
F (x, t) = F (x, k) + (t k) z(F (x, k)) for k < t (k + 1).
Concrete calculations show that this collection of functions has a limit as 0,
that this limit is continuous and satisfies the differential equation (), and also
that any solution of () is a limit of this collection. These calculations give precise information concerning the accuracy of the numerical scheme for computing
approximate solutions described by this approach.
The second proof scheme uses a fixed point theorem. It considers the mapping
F 7 F given by the equation
F (x, t) = x +
15.2
Dynamics on a Manifold
214
Lemma 15.2.3. A composition of two Lipschitz functions is Lipschitz, and a composition of two locally Lipschitz functions is locally Lipschitz.
Proof. Suppose that f : X X is Lipschitz, with Lipschitz constant L, that
(X , d ) is a third metric space, and that g : X X is Lipschitz with Lipschitz
constant M. Then
d (g(f (x)), g(f (y))) Md (f (x), f (y)) LMd(x, y)
for all x, y X, so g f is Lipschitz with Lipschitz constant LM.
Now suppose that f and g are only locally Lipschitz. For any x X there is a
neighborhood U of x such that f |U is Lipschitz and a neighborhood V of f (x) such
that g|V is Lipschitz. Then f |U f 1 (V ) is Lipschitz, and, by continuity, U f 1 (V )
is a neighborhood of x. Thus g f is locally Lipschitz.
In preparation for the next result we note the following immediate consequences
of equation ():
Dh(p)p = h(p)
for all p M and q h(M). We also note that for everything we have done up to
this point it is enough that r 1, but the following result depends on r being at
least 2.
215
216
If W and is a second pair with these properties, then W W satisfies (a), and
uniqueness implies that and agree on W W , so the function on W W
that agrees with on W and with on W satisfies (b). In fact his logic extends
to any, possibly infinite, collection of pairs. Applying it to the collection of all such
pairs shows that there is a maximal W satisfying (a), called the flow domain of
, such that there is a unique : W M satisfying (b), which is called the flow
of . Since the flow agrees, in a neighborhood of any point, with a function derived
(by change of time) from one of those given by Theorem 15.2.5, it is continuous,
and it is C s (1 s r) if is C s .
The vector field is said to be complete if W = M R. When this is the case
each (, t) : M M is a homeomorphism (or C s diffeomorphism is is C s ) with
inverse (, t), and t 7 (, t) is a homomorphism from R (thought of as a group)
to the space of homeomorphisms (or C s diffeomorphisms) between M and itself.
It is important to understand that when is not complete, it is because there
are trajectories that go to in finite time. One way of making this rigorous is
to define the notion of going to as a matter of eventually being outside any
compact set. Suppose that Ip = (a, b), where b < , and C M is compact. If we
had (p, tn ) C for all n, where {tn } is a sequence in (a, b) converging to b, then
after passing to a subsequence we would have (p, tn ) q for some q C, and we
could used the method of the last proof to show that (p, b) W .
15.3
217
Remark: In the theory of dynamical systems we are more interested in the future
than the past. In particular, forward stability is of much greater interest than
backward stability, even though the symmetry t 7 t makes the study of one
equivalent to the study of the other. From this point of view it seems that it would
have been preferable to define the vector field index with (V1) replaced by the
normalization requiring that the vector field x 7 x Tx Rm has index 1.
The remainder of this section is devoted to the proof of Theorem 15.3.2. Fix
V(M) with domain C. The first order of business is to show that can be approximated by a well enough behaved vector field that is defined on a neighborhood
of C.
Since C is compact, it is covered by the interiors of a finite collection K1 , . . . , Kk
of compact sets, with each Ki contained in an open Vi that is the image of a C r
parameterization i . Each i induces an isomorphism between T Vi and Vi Rm ,
so that the Tietze extension theorem implies that there is a vector field on Vi that
agrees with on CVi . There is a partition of unity {i } for K1 . . .K
S k subordinate
to the cover V1 , . . . , Vk , and we may define an extension of to V = i Vi by setting
X
(p) =
i (p)i (p).
pVi
218
2 (v, w) = w,
be the projection onto the second component. We say that p is a regular equilib is nonsingular. (Intuitively, the
rium of if p is an equilibrium of and 2 D (p)
derivative at p of the map q 7 q has rank m.) We need the following local result.
Lemma 15.3.3. Suppose that K V V U Rm with U and V open and
K and V compact, and : U [0, 1] is a C r1 (2 r ) function with
(x) = 1 whenever x K and (x) = 0 whenever x
/ V . Let D be a closed subset
of U, and let f : U Rm be a C r1 function whose zeros in D are all regular.
Then any neighborhood of the origin in Rm contains a y such that all the zeros of
fy : x 7 f (x) + (x)y in D K are regular.
Proof. The equidimensional case of Sards theorem implies that the set of regular
values of f |V is dense, and if y is a regular value of f |V , then all the zeros of fy |K
are regular. If the claim is false, there must be a sequence yn 0 such that for
each n there is a xn V D such that xn is a singular zero of fyn . But V D is
compact, so the sequence {xn } must have a limit point, which is a singular zero of
f |D by continuity, contrary to assumption.
the last result, we first choose a perturbation 1 of 1 such that 1 1 +
PnUsing
i i +
n
X
i i
h=i+1
has only
equilibria in Di1 Ki = (K1 . . . Ki1 ) Ki . At the end of this
P regular
219
find that the index of the equilibrium of the second type is 1, so the vector field
index is indeed uniquely determined by the axioms.
We still need to construct the index. One way to proceed would be to define
the vector field index to be the index of nearby smooth approximations with regular equilibria. This is possible, but the key step, namely showing that different
approximations give the same result, would duplicate work done in earlier chapters.
Instead we will define the vector field index using the characterization in terms of
the fixed point index given in the statement of Theorem 15.3.2, after which the
axioms for the fixed point index will imply (V1)-(V3).
We need the following technical fact.
Lemma 15.3.4. If C M is compact, is a locally Lipschitz vector field defined
on a neighborhood U of C, is the flow of , and (p) 6= 0 for all p C, then there
is > 0 such that (p, t) 6= p for all (p, t) C ((, 0) (0, )).
Proof. We have C = K1 . . . Kk where each Ki is compact and contained in
the domain Wi of a C r parameterization i . It suffices to prove the claim with
C replaced with Ki , so we may assume that C is contained in the image of a C r
parameterization. We can use Lemma 15.2.1 to move the problem to the domain of
the parameterization, so we may assume that U is an open subset of Rm and that
is Lipschitz, say with Lipschitz constant L.
Let V be a neighborhood of C such that V is a compact subset of U, and let
M := max k(p)k and m := min k(p)k.
pV
pV
Therefore the intermediate value theorem implies that |h(p, t) p, vi| M|t| kvk.
Since v may be any unit vector, () follows this.
Now suppose that (p, t) = p for some (p, t) C ((, 0) (0, )). Rolles
theorem implies that there is some s between 0 and t such that
0=
d
h(p, s)
ds
p, pi = h(p,s) , p i,
but among the vectors that are orthogonal to p , the origin is the closest, and
k(p,s) p k Lk(p, s) pk LM|s| < LM < m < kp k,
so this is impossible.
220
We now define the vector field index of the pair (U, ) to be ( (, t)|C ), where
is a nearby C r1 vector field, for all sufficiently small negative t. Since such a
is vector field admissible, the last result (applied to C) implies that (, t)|C is
index admissible for all small negative t, and it also (by Homotopy) implies that
the choice of t does not affect the definition.
We must also show that the choice of does not matter. Certainly there is a
neighborhood such that for 0 and 1 in this neighborhood and all s [0, 1], s =
(0 (, t)|C ) = (1 (, t)|C ).
We now have to verify that our definition satisfies (V1)-(V3). But the result
established in the last paragraph immediately implies (V3). Of course (V2) follows
directly from the Additivity property of the fixed point index. Finally, the flow of
the vector field (x) = x on Rm is (x, t) = et x, so for small negative t there is
an index admissible homotopy between (, t)|Dm and the constant map x 7 0, so
(V1) follows from Continuity and Normalization for the fixed point index.
All that remains of the proof of Theorem 15.3.2 is to show that is locally
Lipschitz and defined in a neighborhood of C, then
ind() = (1)m ( (, t)|C )
for sufficiently small positive t. Since we can approximate with a vector field that
is C r1 and has only regular equilibria, by (V2) it suffices to prove this when C is
a single regular equilibrium. If is one of the two vector fields x 7 x Tx Rm or
x 7 (x1 , x2 , . . . , xm ) Tx Rm on Rm , then (x, t) = (x, t) for all x and t,
so the result follows from the relationship between the index and the determinant
of the derivative.
15.4
Dynamic Stability
221
One of the earliest and most useful tools for understanding stability was introduced by Lyapunov toward the end of the 19th century. A function f : M R is
-differentiable if the -derivative
f (p) =
d
f ((p, t))|t=0
dt
d
d
L((p, t))|t=t = L(((p, t ), t))|t=0 ,
dt
dt
222
15.5
A converse Lyapunov theorem is a result asserting that if a set is asymptotically stable, then there is a Lyapunov function defined on a neighborhood of the
set. The history of converse Lyapunov theorems is sketched by Nadzieja (1990).
Briefly, after several partial results, the problem was completely solved by Wilson
(1969), who showed that one could require the Lyapunov function to be C when
the given manifold is C . Since we do not need such a refined result, we will follow
the simpler treatment given by Nadzieja (1990).
Let M, , W , and be as in the last section. This sections goal is:
Theorem 15.5.1. If A is asymptotically stable, then (after replacing M with a
suitable neighborhood of A) there is a Lyapunov function for A.
The construction requires that the vector field be complete, and that certain
other conditions hold, so we begin by explaining how the desired situation can be
achieved on some neighborhood of A. Let U D(A) be an open neighborhood of
A whose closure (as a subset of Rk ) is contained in M. For any metric on M (e.g.,
the one induced by the inclusion in Rk ) the infimum of the distance from a point
p U to a point in M \ U is a positive continuous function on U, so Proposition
10.2.7 implies that there is a C r function : U (0, ) such that for each p U,
be the graph
1/(p) is less than the distance from p to any point in M \ U. Let M
of :
= { (p, (p)) : p U } U R Rk+1 .
M
are the subsets that are closed in Rk+1 :
The closed subsets of M
223
Since IdM is C r , Lemma 15.2.4 implies that is a locally Lipschitz vector field
Using the chain rule, it is easy to show that
. Let
be the flow of .
on M
((p,
h), t) = (p, t), ((p, t))
for all (p, t) in the flow domain of . Since asymptotic stability is a local property,
let
Z
t
B(p, t) =
(p)), s) ds.
((p,
This has two important consequences. The first is that the speed of a trajectory
(p), t) cannot go
of is never greater than one, so the final component of (p,
to in finite (forward or backward) time. In view of our remarks at the end of
Section 15.2, is complete. The second point is that since is bounded below on
(p), t) : t 0 } is bounded, then (p,
) traverses the
any compact set, if { (p,
entire trajectory of beginning at (p, (p)). It follows that A is asymptotically
Note that if L
is a Lyapunov function for and A,
then it is also a
stable for .
and setting L(p) = L(p,
(p)) gives a Lyapunov
Lyapunov function for and A,
function for |U and A. Therefore it suffices to establish the claim with M and
and .
replaced by M
The upshot of the discussion to this point is as follows. We may assume that
is complete, and that the domain of attraction of A is all of M. We may also
assume that M has a metric d that is completethat is, any Cauchy sequence
convergesso a sequence {pn } that is eventually outside of each compact subset of
M diverges in the sense that d(p, pn ) for any p M.
The next four results are technical preparations for the main argument.
224
If p A, then (p) = 0. If p
/ A, then (p, t)
/ A for all t 0 because t is
invariant, and the last result implies that (p) > 0.
Lemma 15.5.6. is continuous.
Proof. Since (p) d(p, A), is continuous at points in A. Suppose that {pn } is a
sequence converging to a point p
/ A. The last result implies that there are t 0
and tn 0 for each n such that (p) = d((p, t), A) and (pn ) = d((pn , tn ), A).
The continuity of and d gives
lim sup (pn ) lim sup d((pn , t), pn ) = (p).
n
On the other hand d((pn , tn ), A) d(pn , A), so the sequence (pn , tn ) is bounded,
and Lemma 15.5.4 implies that {tn } is bounded below. Passing to a subsequence,
we may suppose that tn t , so that
(p) d((p, t ), A) = lim d((pn , tn ), A) = lim inf (pn ).
n
225
We are now ready for the main construction. Let L : M [0, ) be defined by
Z
L(p) =
((p, s)) exp(s) ds.
0
The rest of the argument verifies that L is, in fact, a Lyapunov function.
Since A is invariant, L(p) = 0 if p A. If p
/ A, then L(p) > 0 because
(p) > 0.
To show that L is continuous at an arbitrary p M we observe that for any
> 0 there is a T such that ((p, T )) < /2. Since is continuous we have
((p , T )) < /2 and |((p , t)) ((p, t))| < /2 for all p in some neighborhood
of p and all t [0, T ], so that
Z T
((p, s)) exp(s) ds <
so that
(p, s) exp(s) ds
exp(s) ds = (p)
because ((p, )) is weakly decreasing with limt ((p, t)) = 0. Therefore L(p) <
0 when p
/ A.
We need one more technical result.
Lemma 15.5.7. If {(pn , tn )} is a sequence such that d(pn , A) and there is a
number T such that tn < T for all n, then d((pn , tn ), A) .
Proof. Suppose not. After passing to a subsequence there is a B > 0 such that
d((pn , tn ), A) < B for all n, so the sequence {(pn , tn )} is contained in a compact
set K. Since the domain of attraction of A is all of M, is continuous, and K
is compact, for any > 0 there is some S such that d((p, t), A) < whenever
p K and t > S. The function p 7 d((p, t), A) is continuous, hence bounded on
the compact set K [T, S], so it is bounded on all of K [T, ). But this is
impossible because tn > T and
d(((pn , tn ), tn ), A) = d(pn , A) .
226
It remains to show that if U is open and contains A, then there is an > 0 such
that L1 ([0, ]) U. The alternative is that there is some sequence {pn } in M \ U
with L(pn ) 0. Since L is continuous and positive on M \ U, the sequence must
eventually be outside any compact set. For each n we can choose tn 1 such that
((pn , 1)) = d((pn , tn ), A), and the last result implies that ((pn , 1)) , so
Z 1
Z 1
L(pn )
((pn , t)) exp(t) dt ((pn , t))
exp(t) dt .
0
This contradiction completes the proof that L is a Lyapunov function, so the proof
of Theorem 15.5.1 is complete.
15.6
This section establishes the relationship between asymptotic stability and the
vector field index. Let M, , and be as before. If A is a compact set of equilibria
for that has a compact index admissible neighborhood C that contains no other
equilibria of , then ind(C ) is the same for all such C; we denote this common
value of the index by ind (A).
Theorem 15.6.1. If A is an ANR that is asymptotically stable, then
ind (A) = (A).
Proof. From the last section we know that (after restricting to some neighborhood
of A) there is a Lyapunov function L for . For some > 0, A = L1 ([0, ]) is
compact. Using the flow, it is not hard to show that A is a retract of A for some
> , and that A is a neighborhood of A , so A is an ANR. For each t > 0,
(, t)|A maps A to itself, and is homotopic to the identity, so
(A ) = ((, t)|A ) = (1)m ind(|A ).
Since A is an ANR, there is a retraction r : C A, where C is a neighborhood
of A. By taking small we may insure that A C, and we may then replace
C with , so we may assume the domain of r is actually A . If i : A C is the
inclusion, then Commutativity gives
(A) = (r i) = (i r) = (r),
so it suffices to show that if t > 0, then
((, t)|C ) = (r).
Let W M M be a neighborhood of the diagonal for which there is convex
combination function c : W [0, 1] M as per Proposition 10.7.9. We claim that if
T is sufficiently large, then there is an index admissible homotopy h : A [0, 1] A
between IdA and r given by
0 t 31 ,
(p, 3tT ),
h(p, t) = c(((p, T ), r((p, T ))), 3(t 13 )), 31 t 23 ,
2
t 1.
r((p, 3(1 t)T )),
3
227
This works because there is some neighborhood U of A such that c((p, r(p)), t) is
defined and in the interior of A for all p U and all 0 t 1, and (A , T ) U
if T is sufficiently large.
The following special case is a prominent result in the theory of dynamical
systems.
Corollary 15.6.2 (Krasnoselski and Zabreiko (1984)). If {p0 } is asymptotically
stable, then
ind ({p0 }) = 1.
Physical equilibrium concepts are usually rest points of explicit dynamical systems, for which the notion of stability is easily understood. For economic models,
dynamic adjustment to equilibrium is a concept that goes back to Walras notion of
tatonnement, but such adjustment is conceptually problematic. If there is gradual
adjustment of prices, or gradual adjustment of mixed strategies, and the agents understand and expect this, then instead of conforming to such dynamics the agents
will exploit and undermine them. For this reason there are, to a rough approximation, no accepted theoretical foundations for a prediction that an economic or
strategic equilibrium is dynamically stable.
Paul Samuelson (1941, 1942, 1947) advocated a correspondence principle, according to which dynamical stability of an equilibrium has implications for the
qualitative properties of the equilibriums comparative statics. Samuelsons writings consider many particular models, but he never formulated the correspondence
principle as a precise and general theorem, and the economics professions understanding of it has languished, being largely restricted to 1-dimensional cases; see
Echenique (2008) for a succinct summary. However, it is possible to pass quickly
from the Krasnoselski-Zabreiko theorem to a general formulation of the correspondence principle, as we now explain.
Let U Rm be open, let P be a space of parameter values that is an open
subset of Rn , and let z : U P Rm be a C 1 function that we understand as
a parameterized vector field. (Working in a Euclidean setting allows us to avoid
discussing differentiation of vector fields on manifolds, which is a very substantial
topic.) For (x, ) U P let x z(x, ) and z(x, ) denote the matrices of partial
derivatives of the components of z with respect to the components of x and
respectively.
We consider a point (x0 , 0 ) with z(x0 , 0 ) = 0 such that x z(x0 , 0 ) is nonsingular. The implicit function implies that there is a neighborhood V of 0 and C 1
function : V U such that (0 ) = x0 and z((), ) = 0 for all V . The
method of comparative statics if to differentiate this equation with respect to ,
using the chain rule, then rearrange, arriving at
d
(0 ) = x z(x0 , 0 )1 z(x0 , 0 ).
d
The last result implies that if {x0 } is asymptotically stable for the vector field
z(, 0 ), then the determinant of x z(x0 , 0 ) is positive, as is the determinant of
d
(0 ) is a positive scalar multiple
its inverse. When m = 1 this says that the vector d
228
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Index
C r , 126
C r -embedding, 144
C r -immersion, 144
C r atlas, 10
C r function, 127
C r manifold, 10, 131
C r submanifold, 11, 136
Q-robust set, 113
Q-robust set
minimal, 114
minimal connected, 114
T1 -space, 66
-limit set, 19, 221
-parameterization, 144
-domination, 17, 104
-homotopy, 17, 104
EXP, 61
FNP, 63
NP, 61
PLS (polynomial local search), 64
PPAD, 64
PPA, 65
PPP (polynomial pigeonhole principle),
64
PSPACE, 61
P, 61
TFNP, 63
Clique, 61
EOTL (end of the line), 64
OEOTL (other end of the line), 65
absolute neighborhood retract, 6, 100
absolute retract, 6, 102
acyclic, 34
affine
combination, 23
dependence, 23
hull, 24
independence, 23
subspace, 24
Alexander horned sphere, 131
algorithm, 60
ambient space, 10, 133
annulus, 144
antipodal function, 203
antipodal points, 200
approximates, 189
Arrow, Kenneth, 2
asymptotic stability, 20, 221
atlas, 10, 131
axiom of choice, 36
balanced set, 116
Banach space, 90
barycenter, 32
base of a topology, 67
bijection, 6
Bing, R. H., 196
Border, Kim, i
Borsuk, Karol, 18
Borsuk-Ulam theorem, 18, 204
bounding hyperplane, 24
Brouwers fixed point theorem, 3
Brouwer, Luitzen, 2
Brown, Robert, i
category, 135
Cauchy sequence, 90
Cauchy-Schwartz inequality, 91
certificate, 61
Church-Turing thesis, 60
closed function, 74
codimension, 24, 136
commutativity configuration, 16, 179
compact-open topology, 83
complete invariant, 194
complete metric space, 90
complete vector field, 216
235
236
completely metrizable, 100
component of a graph, 34
computational problem, 60
complete for a class, 62
computable, 60
decision, 61
search, 61
connected
graph, 34
space, 8, 113, 165
continuous, 78
contractible, 5
contraction, 5
converse Lyapunov theorem, 20, 222
convex, 24
combination, 24
cone, 25
hull, 24
coordinate chart, 10, 131
correspondence, 4, 77
closed valued, 77
compact valued, 4
convex valued, 4, 77
graph of, 77
lower semicontinuous, 78
upper semicontinuous, 77
correspondence principle, 19, 212
critical point, 139, 154
critical value, 139, 154
cycle, 34
Debreu, Gerard, 2
degree, 11, 33, 174
degree admissible
function, 12, 14, 171
homotopy, 12, 171
Dehn, Max, 149
Demichelis, Stefano, 19
derivative, 126, 134, 135
derivative along a vector field, 20
Descartes, Rene, 30
deterministic, 212
diameter, 32
diffeomorphism, 11, 133
diffeomorphism point, 136
differentiable, 126
INDEX
differentiation along a vector field, 221
dimension
of a polyhedron, 26
of a polytopal complex, 30
of an affine subspace, 24
directed graph, 64
discrete set, 131
domain of attraction, 19, 221
domination, 183
dual, 25
Dugundji, James, 93
Dugundji, James, i
edge, 27, 33
Eilenberg, Samuel, 7, 18
Eilenberg-Montgomery theorem, 196
embedding, 6, 131
endpoint, 33, 42
equilibrium, 19, 21
equilibrium of a vector field, 216
regular, 218
essential
fixed point, 7
Nash equilibrium, 8
set of fixed points, 8, 112
set of Nash equilibria, 8
Euclidean neighborhood retract, 6, 99
Euler characteristic, 17, 20, 195
expected payoffs, 37
extension of an index, 183
extraneous solution, 44
extreme point, 29
face, 26
proper, 27
facet, 27
family of sets
locally finite, 85
refinement of, 85
Federer, Herbert, 150
Fermats last theorem, 149
fixed point, 3, 4
fixed point property, 3, 6
flow, 19, 216
flow domain, 216
Fort, M. K., 107
four color theorem, 149
237
INDEX
Freedman, Michael, 149
Fubinis theorem, 150
functor, 135
general position, 41
general linear group, 165
Granas, Andrzej, i
graph, 4, 33
half-space, 24
Hauptvermutung, 196
Hausdorff distance, 70
Hausdorff measre zero, 156
Hausdorff space, 67
have the same orientation, 11, 164
Hawaiian earring, 33, 100
Heegaard, Poul, 149
Hilbert cube, 93
Hilbert space, 91
homology, 2, 3, 177, 196
homotopy, 5, 5859
class, 5
extension property, 198
invariant, 18, 197
principle, 178
homotopy extension property, 103
Hopfs theorem, 18, 197198
Hopf, Heinz, 18, 196
hyperplane, 24
identity component, 165
immersion, 138
immersion point, 136
implicit function theorem, 127
index, 9, 15, 16, 177, 179
index admissible
correspondence, 15, 177
homotopy, 178
vector field, 20
index base, 15, 177
index scope, 15, 178
inessential fixed point, 7
initial point, 27
injection, 6
inner product, 91
inner product space, 91
invariance of domain, 18, 207
invariant, 19
invariant set, 221
inverse function theorem, 127
isometry, 52
Kakutani, Shizuo, 4
Kinoshita, Shinichi, 6, 95, 108
labelling, 51
Lefschetz fixed point theorem, 17, 196
Lefschetz number, 17, 196
Lefschetz, Solomon, 17, 196
Lemke-Howson algorithm, 3649, 62, 64
Lesbesgue measure, 150
lineality space, 26
linear complementarity problem, 45
Lipshitz, 212
local diffeomorphism, 138
locally C r , 130
locally closed set, 98
locally Lipschitz, 212
locally path connected space, 102, 142
lower semicontinuous, 78
Lyapunov function, 221
Lyapunov function for A M, 20
Lyapunov theorem, 221
Lyapunov, Aleksandr, 20
manifold, 10
C r , 131
manifold with boundary, 13, 144
Mas-Colell, Andreu, 17, 115
maximal, 34
measure theory, 150
measure zero, 151, 157
mesh, 32
Milnor, John, 150, 196
Minkowski sum, 111
Moise, Edwin E., 196
Montgomery, Deane, 7, 18
Morse-Sard theorem, 157
moving frame, 165
multiplicative, 16, 180
Mobius, August Ferdinand, 149
narrowing of focus, 183
Nash equilibrium
accessible, 43
238
mixed, 37
pure, 37
refinements of, 8
Nash, John, 2
negatively oriented, 164, 168
negatively oriented relative to P , 169
neighborhood retract, 98
neighbors, 33
nerve of an open cover, 105
no retraction theorem, 100
norm, 90
normal bundle, 140
normal space, 67
normal vector, 24
normed space, 90
opposite orientation, 164, 168
oracle, 62
order of differentiability, 126
ordered basis, 164
orientable, 11, 168
orientation, 163171
orientation preserving, 11, 53, 169
orientation reversing, 11, 53, 169
orientation reversing loop, 167
oriented -manifold, 168
oriented intersection number, 169
oriented manifold, 11
oriented vector space, 11, 164
paracompact space, 85
parameterization, 10, 131
partition of unity, 86
C r , 128
path, 33, 165
path connected space, 164
payoff functions, 37
Perelman, Grigori, 149
Picard-Lindelof theorem, 212, 215
pivot, 57
pivoting, 48
Poincare conjecture, 149
Poincare, Henri, 149
pointed cone, 26
pointed map, 113
pointed space, 113
polyhedral complex, 30
INDEX
polyhedral subdivision, 30
polyhedron, 26
minimal representation of, 27
standard representation of, 27
polytopal complex, 30
polytopal subdivision, 30
polytope, 29
simple, 46
positively oriented, 11, 164, 168
positively oriented relative to P , 169
predictor-corrector method, 58
prime factorization, 63
quadruple
edge, 42
qualified, 42
vertex, 42
quotient topology, 80
Rado, Tibor, 149, 196
recession cone, 25
reduction, 62
regular fixed point, 9
regular point, 11, 139
regular space, 67
regular value, 11, 139
retract, 6, 97
retraction, 97
Ritzberger, Klaus, i, 19
Samuelson, Paul, 19
Sards theorem, 182, 218
Scarf algorithm, 56
Scarf, Herbert, 18
separable, 6
separable metric space, 91
separating hyperplane theorem, 24
set valued mapping, 4
simplex, 31
accessible completely labelled, 58
almost completely labelled, 54
completely labelled, 52
simplicial complex, 31
abstract, 32
canonical realization, 32
simplicial subdivision, 31
simply connected, 149
239
INDEX
slack variables, 44
slice of a set, 153
Smale, Stephen, 149
smooth, 11, 156
Sperner labelling, 51
star-shaped, 5
Steinitz, Ernst, 196
step size, 58
Sternberg, Shlomo, 150
strategy
mixed, 37
pure, 37
totally mixed, 39
strategy profile
mixed, 37
pure, 37
strong topology, 83
strong upper topology, 78
subbase of a topology, 67
subcomplex, 30
submanifold, 11
neat, 146
submersion, 138
submersion point, 136
subsumes, 183
support of a mixed strategy, 63
surjection, 6
tableau, 47
tangent bundle, 133
tangent space, 10, 133
tatonnement, 18
Tietze, Heinrich, 196
topological space, 66
topological vector space, 88
locally convex, 89
torus, 10
trajectory, 212, 213
transition function, 10
translation invariant topology, 88
transversal, 139, 146, 158
tree, 34
triangulation, 31
tubular neighborhood theorem, 140
Turing machine, 59
two person game, 37
Ulam, Stanislaw, 18
uniformly locally contractible metric space,
101
upper semicontinuous, 4, 77
Urysohns lemma, 87
van Dyke, Walther, 149
vector bundle, 140
vector field, 19, 158, 213
along a curve, 165
index admissible, 216
vector field homotopy, 217
index admissible, 217
vector field index, 216
vertex, 27, 32
vertices, 33
connected, 34
Vietoris topology, 68
Vietoris, Leopold, 66
von Neumann, John, 4
Voronoi diagram, 30
walk, 33
weak topology, 83
weak upper topology, 80
well ordering, 85
well ordering theorem, 85
Whitney embedding theorems, 133
Whitney, Hassler, 10
wild embedding, 131
witness, 61
zero section, 140, 158