EDUCATIVE COMMENTARY ON

KVPY 2018 MATHEMATICS QUESTIONS

Contents

Section 1 3

Section 2 25

Concluding Remarks 46

Introduction

I have been writing educative commentaries on the Mathematics sections

of the JEE (Advanced) papers since 2003. These have been supplemented by
similar commentaries on ISC Entrance tests from 2014 onwards. Some read-
ers have suggested that I should also write commentaries on other competitive
annual examinations such as the RMO’s (Regional Mathematics Olympiads),
the INMO (Indian National Mathematics Olympiads), the IMO (Interna-
tional Mathematics Olympiads), the CMI Entrance tests and so on. Clearly,
this is a very time consuming task. Moreover, already there are websites
which give solutions to the problems in many of these examinations.
But the KVPY (Kishore Vaigyanik Protsahan Yojana) stands out. Unlike
JEE, it was not introduced as an entrance test at its inception in 1999,
although its syllabus is comparable to it. Rather it is meant to encourage the
HSC students to take up careers in sciences. It is not a test taken by a huge
number of candidates. In fact, I had not heard much about it. Occasionally,
someone would ask me an interesting question without revealing that it was
from KVPY. This year someone asked me a question involving asymptotic
estimation of an integral and later revealed that it was from KVPY 2018. (It
is Q.86 in this commentary.) I found there is considerable educative content
in it. Then I looked at all the mathematics questions of KVPY 2018. I found
that the questions in Section I (numbering 1 to 20) are comparable to those
in JEE. But some of those in Section 2 (numbering 81 to 90) are superior in
terms of their variation and the ideas they inspire (which is the true worth
of a good problem). In the past the JEE also had a few such questions. But
nowadays, they are rare.
The format of the KVPY questions is simple, with only one correct answer
to each question. True, this sometimes makes it possible to guess the correct
answer by elimination. But such cases are few. And, in some of them, it is
probably intentional.
So, as an experiment, I am presenting an educative commentary on KVPY
2018. The style and format resemble the commentaries on JEE papers. How-
ever, when a particular question is a routine computational question, the
solution will be given tersely. Such questions appear mostly in Section 1. In
Section 2, there are a few questions where an intuitive reasoning is possible
and perhaps intended. It will be given as a justification in some cases.
Unless otherwise stated, all the references made are to the author’s book
Educative JEE (Mathematics) published by Universities Press, Hyderabad.
The third edition of this book is now available in the market.
Readers who notice any errors in this commentary are invited to send an
email to the author at kdjoshi314@gmail.com or send an SMS or a WhatsApp
message to the author at 9819961036. Alternate solutions and any other
comments are also solicited.
I am indebted to Prof. Srinivas Sastry who asked me Q.86, which got
me interested in the KVPY, to Prof. Gopal Srinivasan for the asymptotic
estimation in it and to Siddhesh Naik for the solution to Q.90 and also
for drawing my attention to the surprisingly simple formula for the area
subtended by a chord of a parabola in the solution to Q.89.

K. D. Joshi

2
 Section 1
" #
a b
Q.1 Suppose A = is a real matrix with non-zero entries, ad−bc 6= 0
c d
and A2 = A. Then a + d equals

A. 1 B. 2 C. 3 D. 4

Answer and Comments: (A). By a direct calculation, the entry in

the first row and second column of A2 is ab + bd. This equals b since
A2 = A. As b 6= 0, we must have a + d = 1.
Since all the entries are given to be non-zero, we got an easy solution
by equating only one entry of A2 with the corresponding entry of A. It
would thus appear that the hypothesis ad − bc 6= 0 is redundant. The
hypothesis is needed when some of the entries of A may vanish. For
example, when A = I2 , the identity matrix, we have A2 = A. But then
a + d equals 2.
There is a more sophisticated solution using what are called eigen-
values of a matrix. A number λ is called an eigenvalue of an n × n
matrix A if there is some non-zero column vector u such that Au = λu.
In general an n × n matrix A has n (not necessarily distinct) eigenval-
ues, some of which may be complex. The product and the sum of these
eigenvalues equal, respectively, the determinant and the trace (i.e. the
sum of the diagonal elements) of A. So, the present problem asks us to
calculate the trace of A, given that its determinant is 0, which means
that at least one eigenvalue is 0.
It is easy to show that if λis an eigenvalue of A, then λ2 is an
eigenvalue of A2 . Hence from A2 = A we get that the only possible
eigenvalues of A are 1 and 0. Therefore the trace is the number of
eigenvalues that equal 1. (If that sounds awkward, take it to mean the
multiplicity of the eigenvalue 1, defined analogously to the multiplicity
of the root of a polynomial. But we shall not go into it.)
In our problem, n = 2 and at least one of the eigenvalues is 0. But
both cannot be 0, as otherwise A would be the zero matrix. Therefore
one eigenvalue is 0 and the other 1. Hence the trace, which equals their
sum is 1.

3
 Incidentally, a matrix A which satisfies A2 = A is called an
idempotent matrix because all its powers, A, A2 , A3 , . . . , are the same.

Q.2 On any given interval of positive length on the unit circle |z| = 1 in the
complex plane,

A. there need not be any root of unity

B. there lies exactly one root of unity
C. there are more than one but finitely many roots of unity
D. there are infinitely many roots of unity

Answer and Comments: (D). It would have been better had the
question included a definition of a root of unity, not that it is hard
to guess, once you know that ‘unity’ is an old name for the number
1 (‘naught’ being an old name for the number 0). So, logically, a
complex number z is a root of unity if there is some positive integer
n such that z n = 1. For every positive integer n, there are precisely
n distinct n-th roots of unity. They are the complex numbers of the
form zk = e2πki/n where k = 0, 1, 2, . . . , n − 1. (Other values of k do
not give new roots because e2πri/n = e2πsi/n whenever r and s differ by
an integral multiple of n. See Chapter 7 for more on complex roots of
unity and their applications in trigonometry.)
p
Put differently, if r = is a rational number where p, q are positive
q
integers, then z = e2πri is a root of unity, specifically, a q-th root
of unity. An arc of the unit circle will be a set of the form C =
{z ∈ C | : z = eiθ , α ≤ θ ≤ β} where α, β are some real numbers
with 0 ≤ α < β < 2π. It will contain the root e2πpi/q if and only if
α ≤ 2πp q
α
≤ β, which is equivalent to 2π ≤ pq ≤ 2π
β
. So there are as many
roots of unity in the arc C as there are rational numbers in the interval
α β
[ 2π , 2π ].
Every open interval contains infinitely many rationals. Moreover,
since β−α
2π
α β
< 1, no two distinct rationals in the interval [ 2π , 2π ] differ by
an integer. So the corresponding roots of unity are distinct. Hence (D)
is true.
A simple problem once you get the correct thought.

4
Q.3 For 0 < θ < π2 , four tangents are drawn at four points (±3 cos θ, ±2 sin θ)
2 2
to the ellipse x9 + y4 = 1. If A(θ) denotes the area of the quadrilateral
formed by these four tangents, the minimum value of A(θ) is

A. 21 B. 24 C. 27 D. 30

Answer and Comments: (B). An absolutely straight forward prob-

lem. By symmetry, the quadrilateral will be a rhombus with diagonals
along the coordinate axes. The portion in the first quadrant is the
triangle OAB where A and B are the points where the tangent at the
point (3 cos θ, 2 sin θ) cuts the axes. The equation of the tangent is
x cos θ
3
+ y sin
2
θ
= 1. Hence A(θ) = 4 × sin θ3cos θ = sin242θ which is minimum
when sin 2θ = 1.
√
Q.4 Let S = {x ∈ IR : cos(x) + cos( 2x) < 2}. Then

A. S = ∅
B. S is a non-empty finite set
C. S is an infinite proper subset of IR\{0}
D. S = IR\{0}

Answer and Comments: √ (D). Clearly 0 ∈ / S since for x = 0, the

expression cos(x) + cos( 2x) attains√its maximum value 2. For any
other x ∈ IR, we have cos(x)√+ cos( 2x) ≤ 2 with equality possible
only √when both cos x and cos( 2x) equal 1 each. That means x = 2mπ
and 2x = 2nπ for some non-zero integers
√ √ m, n. But that would mean
2nπ n
2 = 2mπ = m . This contradicts that 2 is irrational. So every x 6= 0
is in S.
Another simple problem once you get √the correct throught, viz. that
2 is the maximum value of cos x + cos( 2x) and is attained only√when
both the terms attain their maximum values. We could replace 2 by
any irrational number. That would make the problem
√ more general but
also easier to solve. By giving the irrational as 2, there is a temptation
to relate it to the cosine function by the equation cos( π4 ) = √12 . But
√
that is useless. All that is needed is that 2 is irrational.

5
Q.5 On a rectangular hyperbola x2 − y 2 = a2 , a > 0, three points A, B, C
are taken as follows: A = (−a, 0); B and C are placed symmetrically
w.r.t. the x-axis on the branch of the hyperbola not containing A.
Suppose that the triangle ABC is equilateral. If the side length of the
triangle ABC is ka, then k lies in the interval

A. (0, 2] B. (2, 4] C. (4, 6] D. (6, 8]

Answer and Comments: (B). Clearly the answer is independent of

a. So, by taking a as the unit of length, we may suppose a = 1. (This
is not much of a simplification, but spares you having to write the
redundant a in every step.)
So we can take A = (−1, 0), B = (sec θ, − tan θ) and C = (sec θ, tan θ)
for some θ ∈ [0, π/2). qThen the equilaterality of the triangle gives
2 tan θ = BC = CA = (sec θ + 1)2 + tan2 θ. Squaring and simplify-
ing, this reduces to tan2 θ = 1+sec θ and further to sec2 θ−1 = 1+sec θ.
Since 1 + sec θ never vanishes, we get sec θ −
√ 1 = 1, i.e. sec θ√= 2. The
side √of the triangle, k, equals 2 tan θ = 2 sec2 θ − 1 = 2 3. Since
1 < 3 < 2, k lies in (2, 4).

Q.6 The number of real solutions x of the equation

1
cos2 (x sin(2x)) + = cos2 x + sec2 x
1 + x2
is

A. 0 B. 1 C. 2 D. infinite

Answer and Comments: (B). An optimisation problem converted

artificially to a problem of solutions of an equation. (See Chapter 14
for more such examples.) By inspection, x = 0 is a solution. By the
A.M.-G.M. inequality, the R.H.S. is at least 2 while every term on the
L.H.S. is at most 1. So the only way equality can hold is when every
term on both the sides equals 1. And that happens only for x = 1.
In Q.4 too the key idea was that for equality to hold in cos(x) +
√
cos( 2x) = 2, every term on the L.H.S. must attain its maximum. Two
problems in the same paper based on the same trick are too many.

6
 2 2
Q.7 Let xa2 + yb2 = 1, a > b, be an ellipse with foci F1 , F2 . Let AO be
its semi-minor axis, where O is the centre of the ellipse. The lines
AF1 and AF2 , when extended, cut the ellipse again at points B and C
respectively. Suppose that the triangle ABC is equilateral. Then the
eccentricity of the ellipse is

A. √1 B. √1 C. 1
D. 1
2 3 3 2

Answer and Comments: (D). Yet another straightforward computa-

tional problem about conics. Like Q.5, here too we have an equilateral
triangle. We take F1 = (ae, 0), F2 = (−ae, 0) and A as (0, −b). Then
B will be of the form (a cos θ, b sin θ) for some θ ∈ (0, π2 ). Collinearity
of A, F1 and B gives, on simplification,
cos θ
e= (1)
1 + sin θ
By symmetry about the minor axis, C = (−a cos θ, b sin θ). Equating
(BC)2 with (AB)2 and using that b2 = a2 (1 − e2 ) we get

3 cos2 θ = (1 − e2 )(1 + sin θ)2 (2)

Thus we have a system of two equations in two unknowns, viz. e and

θ. We could determine both. But as our interest is only in e, we try
to eliminate θ between these two equations. This is very easy in the
present problem and gives 3e2 = 1 − e2 from which e = 12 .

Q.8 Let a = cos 1◦ and b = sin 1◦ . We say that a real number is algebraic if
it is a root of a polynomial with integer coefficients. Then

A. a is algebraic but b is not algebraic

B. b is algebraic but a is not algebraic
C. both a and b are algebraic
D. neither a nor b is algebraic.

Answer and Comments: (C). It is to be noted that the problem does

not ask to identify the polynomials, if any, with integer coefficients that
have a and b as roots. We merely have to see if such polynomials exist.

7
Expressions for the sines and cosines of angles of the form k degrees
are available when k = 180, 90, 60, 45, 36, 30, 15 and a few other values,
but not when k = 1. But if we recall how some of these values are
obtained, we see that they crucially use the formulas for the sines and
cosines of an angle kα in terms of the sines and cosines of α. For
example, cos 30◦ = 2 cos2 15◦ − 1. So, if we know cos 30◦ , we get cos 15◦
as a root of a quadratic.
In general for every integer n ≥ 0, cos(2n + 1)θ and sin(2n + 1)θ can
be expressed as polynomials with integer coefficients of degree 2n + 1
in cos θ and sin θ respectively. The cases n = 0, 1 are very familiar, viz.
cos 3θ = 4 cos3 θ − 3 cos θ and sin 3θ = 3 sin θ − 4 sin3 θ. The general
case can be proved by induction on n using the formulas

cos(2n + 1)θ = cos(2n − 1)θ cos 2θ − sin(2n − 1)θ sin 2θ (1)

sin(2n + 1)θ = sin(2n − 1)θ cos 2θ + cos(2n − 1)θ sin 2θ (2)

but the proof gets rather complicated because in the inductive step we
also need that cos(2n − 1)θ and sin(2n − 1)θ contain only odd degree
terms in cos θ and sin θ respectively. (So, this is an example where it
is easier to prove by induction, a more specific statement.)
Anyway, if we assume this result, then we can express cos 45◦ as a
polynomial of degree 45 in a (i.e. cos 1◦ ), say
1
√ = c45 a45 + c43 a43 + . . . + c3 a3 + c1 a (3)
2
where c45 , c43 , . . . , c3 and c1 are some integers. Squaring and multiply-
ing by 2, we get

2(c45 a45 + c43 a43 + . . . + c3 a3 + c1 a)2 = 1 (4)

which shows that a is algebraic because all coefficients are integers.

Similarly, by expressing sin 45◦ as a polynomial with integer coefficients
in b, we can show that b is algebraic.
This is a rather complicated argument and will take a long time
to write out fully. But in a multiple choice question, once a candidate
realises that (3) and a similar equation for b is possible, even without
proving it he can answer that both a, b are algebraic.

8
 Interestingly, if we consider the complex number z = a + ib =
cos 1 +i sin 1◦ , then it is much easier to show that z is algebraic because
◦

by DeMoivre’s rule, z 360 = cos 360◦ +i sin 360◦ = 1. More generally, this
argument shows that any root of unity is an algebraic complex number.
However, deducing from this that its real and imaginary parts are both
algebraic is not so easy. Clearly, z = a − ib is also algebraic because
(z)n = z n . It takes some work to show that the sum of two algebraic
numbers is algebraic. Once we grant that, we get that 2a = z + z and
2ib = z − z are algebraic. From there, the algebraicity of a and b is but
one step away.
There is, in fact, a way using DeMoivre’s rule, to express cos nθ
and sin nθ as certain polynomials in cos θ and sin θ. The resulting
polynomials are called Chebyshev or Tchebyshev polynomials. See
Exercise (7.21), parts (d) and (e).
Q.9 A rectangle with its sides parallel to the x-axis and the y-axis is in-
scribed in the region bounded by the curves y = x2 −4 and 2y = 4 −x2 .
The maximum possible area of such a rectangle is closest to the integer

A. 10 B. 9 C. 8 D. 7

Answer and Comments: (B). The region bounded by the two parabo-
2
las y = x2 −4 and y = 4−x
2
and an inscribed rectangle P QRS are shown
in the figure below.
y

2 4
y=x−
(0,2)

R Q

x
(−2,0) O (2,0)

S P

Clearly the rectangle is symmetric about the y-axis. Assume that its
2
base is 2a. Then P = (a, a2 − 4) and Q = (a, 4−a
2
). Hence the area of

9
 2
the rectangle is 2a × ( 4−a
2
− (a2 − 4)) which comes out as a(12 − 3a2 ).
Call this as f (a). We have to maximise this as a varies over [0, 2].
Setting f ′ (a) = 0 gives a = √23 . As this is the only critical point and
f (0) = f (2) = 0, f attains its maximum at a = √23 . The maximum
q √
area is √23 × (12 − 4) = √163 = 256 3
= 85.33... Among the given
options the one whose square is closest to 85.33.. is 9.

Q.10 Let f (x) = x| sin x|, x ∈ IR. Then

A. f is differentiable for all x, except at x = nπ, n = 1, 2, 3, . . .

B. f is differentiable for all x, except at x = nπ, n = ±1, ±2, ±3, . . .
C. f is differentiable for all x, except at x = nπ, n = 0, 1, 2, 3, . . .
D. f is differentiable for all x, except at x = nπ, n = 0, ±1, ±2, ±3, . . .

Answer and Comments: (B). The functions x and sin x are dif-
ferentiable everywhere. But |x| is not differentiable at x = 0. So,
| sin x| is not differentiable when sin x = 0, i.e. when x = nπ for
n = 0, ±1, ±2, ±3, . . .. For n 6= 0, the factor x is non-zero and hence
the product x| sin x| is non-differentiable at x = nπ. However, when
n = 0, we have to check the differentiability of f (x) from the first prin-
ciples. In a small neighbourhood N of 0, we have f (x) = x sin x for
x > 0 and f (x) = −x sin x for x < 0. Then the right handed derivative
x sin x − 0
f+′ (0) equals lim+ = lim+ sin x = 0. Also the left handed
x→0 x x→0
′ −x sin x − 0
derivative f− (0) equals lim− = lim− − sin x = 0. As both
x→0 x x→0
these are equal, f is differentiable at 0.

Q.11 Let f : [−1, 1] −→ IR be a function defined by

(  
x2 cos( πx ) for x 6= 0
 
f (x) =
0 for x = 0

The set of points where f is not differentiable is

A. {x ∈ [−1, 1] : x 6= 0}
2
B. {x ∈ [−1, 1] : x = 0 or x = n+1
,n ∈ ZZ}

10
 2
C. {x ∈ [−1, 1] : x = n+1
,n ∈ ZZ}
D. [−1, 1]

Answer and Comments: (C). This is very similar to the last ques-
tion. But the first factor is x2 and we have to check if it is strong
 enough

to cure the non-differentiability of the second factor, viz. cos( πx ). It
 

is differentiable except when cos( πx ) = 0, i.e. except when πx is an odd

multiple of π2 , which is equivalent to saying that x is of the form 2n+1
2
for
2
n ∈ IN. Near  all these
 points, the first factor x is non-zero and so the
2 π 
product x cos( x ) will remain non-differentiable. However, in a neigh-
bourhood of x = 0, the second factor lies in [0, 1] and so f (x) ∈ [0, x2 ]
f (x) − f (0)
for x 6= 0. So, lies between 0 and x for x > 0 and between
x
−x and 0 for x < 0. In both the cases, it tends to 0 as x → 0. So, f is
differentiable at 0, with f ′ (0) = 0. Hence (C) is correct.
In the last problem, the differentiability of the product x| sin x|
at x = 0, despite the non-differentiability of the second factor was
possible because near 0, sin x is comparable to x and so x| sin x| is
comparable to ±x2 which is differentiable at 0 regardless of which sign
holds on which side  of 0.  In the present problem, the differentiability
of the product x cos( πx ) at x = 0, despite the discontinuity of the
2

second factor at 0 was possible because the second factor was bounded
in a neighbourhood of 0. The moral is that when you are dealing
with a product of two functions, even if one of the factors is bad, the
other factor is sometimes strong enough to cure this defect. This does
not apply to sums. If h(x) = f (x) + g(x) and g is non-differentiable
at some point a, then differentiability of f at a always means non-
differentiability of h at a because otherwise we would have g = h − f
differentiable at a. So, the goodness of one term does not cure the
badness of the other. Interestingly, sometimes the sum of two bad
functions can be good. An extreme example is when g = −f . Here
the defects of the two terms cancel each other and so together they
live happily like a homely looking girl married to a hardly working boy.
(The descriptions of the qualifications of both the parties are borrowed
from popular matrimonial advertisements in cheap newspapers!)
Problems dealing with such good-bad pairs, especially their products,
have been asked many times in the JEE. It is important to note that if

11
 we are dealing with some h(x) = f (x)g(x) where f (x) is good in some
sense, say differentiability at some point a, and g(x) is bad, then for the
product h(x) to be good despite the badness of g(x) at a, f (a) must
be 0. For otherwise, by continuity, f (x) would be non-zero in some
neighbourhood N of a and in this neighbourhood we can write g(x) as
h(x)
f (x)
which will make it good, being the ratio of two good functions.
Rπ
Q.12 The value of the integral (1 − | sin 8x|)dx is
0

A. 0 B. π − 1 C. π − 2 D. π − 3

Answer and Comments: (C). The term | sin 8x| changes its def-
inition whenever x crosses a point at which sin 8x vanishes. In the
interval [0, π], sin 8x = 0 when x = kπ 8
for k = 0, 1, . . . , 8. So we
π
consider 8 intervals of the form [(k − 1) 8 , k π8 ] for k = 1, 2, . . . , 8. Cor-
respondingly, the given integral splits into 8 integrals, one over each
one of these subintervals. But it is hardly necessary to evaluate each
of these 8 integrals separately. Because of the periodicity of the sine
function, sin(8( π4 + x)) = sin 8x and so the graphs of f (x) on alternate
intervals will be the same, while on two adjacent intervals of the form
[(k−1)π/8, kπ/8] and [kπ/8, (k+1)π/8], where k is odd, the graphs will
be reflections of each other in the line x = kπ/8. Hence the integral
over each subinterval will be the same. Therefore the given integral
π/8
equals 8 times the integral (1 − sin 8x)dx which comes out to be
R
0
π/8
π
+ 81 cos 8x = π
− 14 . Therefore the given integral equals π − 2.

8 0 8

An imaginatively designed problem based on the periodicity of

the sine function. Those who notice the intended short cut will be
rewarded. That is a characteristic of many good problems. They can
be done by brute force, like evaluating all eight integrals in the present
problem. But the time taken will be atrocious and that itself prompts
a good candidate to see if there is some short cut.

Q.13 Let ln x denote the logarithm of x with respect to the base e. Let S ⊂ IR
be the set of all points where the function ln(x2 − 1) is well defined.
Then the number of functions f : S −→ IR that are differentiable,
satisfy f ′ (x) = ln(x2 − 1) for all x ∈ S and f (2) = 0 is

12
 A. 0 B. 1 C. 2 D. infinite

Answer and Comments: (D). A differential equation of order 1

is given about the function f (x) along with an initial value f (2) =
0. We are taught that there is a unique solution to this initial value
problem. But the catch is that in all such theorems, there is an inherent
assumption that the domain of the function is a connected set, i.e. a
single interval (which may be finite or infinite). In the present problem,
the set S, which is the domain of f , consists of all x for which x2 > 1.
This is a union of two disjoint open intervals (−∞, −1) and (1, ∞).
The point 2 lies in the second interval and on it there will be only
one solution. But the value of f at no point of (−∞, −1) is specified
or can be derived from the solution of the d.e. on the other interval.
The two intervals are like two separated islands in a sea. Little can
be said about what happens on one of them from what happens on
the other. By taking f (−2) arbitrarily, there will be infinitely many
possible functions which satisfy all the conditions of the problem.
A very good problem, requiring hardly any computation, but testing
conceptual perfection, in this case an awareness that the domain must
be an interval. So many familiar theorems such as the Mean Value
Theorems or the Intermediate Value Property for continuous functions
break down when they are applied to a function whose domain is the
the disjoint union of two (or more) open intervals.

Q.14 Let S be the set of all real numbers p such that there is no non-zero
Rx
continuous function f : IR −→ IR satisfying f (t)dt = pf (x) for all
0
x ∈ IR. Then S is

A. the empty set

B. the set of all rational numbers
C. the set of all irrational numbers
D. the whole IR

Answer and Comments: (D). The formulation of the question is un-

necessarily clumsy. It would have been more natural to ask to identify

13
the set of those real numbers p for which at least one function f (x) sat-
isfying the conditions given exists. That set would be the complement
of the set S in the question and so the problem would be the same.
Instead of asking every man for whom there is no woman who can be
called his wife to raise his hand, it is far better to say ‘Those who are
married may raise their hands’ !
Anyway, coming to the question, in the equation
Z x
f (t)dt = pf (x) (1)
0

the L.H.S. is a function defined by an integral and by the Fundamen-

tal Theorem of Calculus (second form) it is differentiable (because f
is given to be continuous) and moreover, its derivative is f (x). So,
differentiating both the sides of (1),

f (x) = pf ′ (x) (2)

for all x ∈ IR. Note that p 6= 0 as otherwise f (x) would vanish identi-
cally. So we assume p 6= 0. Then (2) is a differential equation whose
general solution is

f (x) = kex/p (3)

where k is an arbitrary constant. Note that (1) implies that f (0) = 0.

So, putting x = 0 in (3), we get

0 = ke0 = k (4)

which makes f (x) vanish identically, a contradiction. Hence there is

no value of p for which there exists some continuous function which
satisfies (1) for all x. Therefore (D) is the correct option.
A clumsily worded question, lacking a clear purpose. One possible
conclusion that can be drawn is that a function obtained by integrating
some other function cannot be a constant multiple of that function.
It generally grows much faster at least when the integrand is positive.
Note also that if f (x) is a polynomial of degree n with a positive leading
Rx
coefficient, then f (t)dt is a polynomial of degree n + 1 which grows
0
faster than f (x).

14
Q.15 The probability of men getting a certain disease is 12 and that of women
getting the same disease is 51 . The blood test that identifies the disease
gives the correct result with probability 45 . Suppose a person is chosen
at random from a group of 30 males and 20 females, and the blood test
of that person is found to be positive. What is the probability that the
chosen person is a man?
75 3 15 3
A. 107
B. 5
C. 19
D. 10

Answer and Comments: (A). The stipulation about the probability

of the blood test giving the correct result has two possible meanings. In
real life, it means that if a person tests positive then his/her chances of
having that disease are 4/5. To dispel uncertainty, more sophisticated
tests such as biopsies are performed. But in this problem, the term
is used with a different meaning. Suppose that the test is performed
on a person x. Let p1 be the probability that x tests positive when x
has the disease and p2 the probability that x tests positive even when
x does not have the disease. In general p1 , p2 may not be complemen-
tary probabilities. (This happens, when the test, e.g. high fever, may
come out to be positive for a variety of reasons, not all related to the
particular disease.) In the problem given, however, the meaning of the
sentence about correct identification of the disease is apparently that
p1 = 45 and p2 = 15 . That is, a person having the disease will test
positive with probability 54 and a healthy person will test positive with
probability 15 .
Assuming that this is the intended interpretation, the desired
probability, say p, is the conditional probability that a person tested is
a man, given that the test is positive. So, let us first find the probability
P (E) where E is the event that the test is positive. E is the disjunction
of two mutually exclusive events E1 and E2 , where E1 is the event that
the person is diseased and the test is correct and E2 is the event that
the person is healthy but the test is incorrect. So

P (E) = P (E1 ) + P (E2 ) (1)

To calculate these two probabilities, let S be the set of all 50 persons,

M the subset of men and W the subset of women in it. Let D be the
subset of diseased persons and H the set of healthy persons. The data

15
of the problem means that |M| = 30, |W | = 20, |M ∩ D| = 12 |M| = 15
and |W ∩ D| = 51 |W | = 4. It is convenient to show this in a Venn
diagram, where the set D of diseased persons is shown in the middle
and the figures in the various regions represent their cardinalities. (Such
diagrams are only indicative and so it is not necessary that the figures
be proportional to the areas, as is the case in a pie cut diagram.)

4
15 16
D

W
15

Now, for E1 to occur, the person must be in D and the test correct.
The probabilities of these two events are 19
50
and 45 respectively. As
the correctness of the test does not depend on whether the person is
diseased or not, these two sub-events are independent of each other.
Therefore,
19 4 76
P (E1 ) = × = (2)
50 5 250
Similarly, for E2 to occur, the person must be healthy but the test
wrong. So
31 1 31
P (E2 ) = × = (3)
50 5 250
Putting these into (1), we get
76 31 107
P (E) = + = (4)
250 250 250
So we have found the probability that a randomly chosen person tests
positive. Now, given this, we have to find the probability that the

16
person tested is a man. For this we need to find P (E ∩ M), i.e. the
probability that the person chosen is a man and tests positive. Once we
do that, Bayes theorem will give that the ratio P (E∩M
P (E)
)
is the probability
that a person testing positive is a man.
Clearly,

P (E ∩ M) = P (M)P (F/M) (5)

where P (M) is the probability that the chosen person is a man and
P (F/M) is the probability that a man chosen at random tests positive.
As there are 30 men in a population of 50, P (M) is easy to find. It
comes as
30 3
P (M) = = (6)
50 5
But finding P (F/M), i.e. the probability that a randomaly chosen
man tests positive needs some work. The caclulations are very similar
to those for P (E), except that we replace the entire population set S
with the set M of all men. Then M ∩ D is the set of diseased men and
has 15 elements in it while M ∩ H is the set of healthy men and also
has 15 elements in it. So
15 4 15 1 75 1
P (F/M) = × + × = = (7)
30 5 30 5 150 2
From (5), (6) and (7),
1 3 3
P (E ∩ M) = × = (8)
2 5 10
So, finally by Bayes’ theorem, the desired probability p, i.e. the prob-
ability that a person testing positive is a man is

P (E ∩ M) 3/10 75
p = P (M/E) = = = (9)
P (E) 107/250 107

Probability problems often have a real life setting which makes them
very appealing. But such problems are also often prone to controversies
arising from different interpretations of the data. For example, in this

17
problem, let us take the first interpretation that if a person x tests pos-
itive then he/she is in D with probability 45 and in H with probability
1
5
. If x ∈ D, then x is a man with probability 15 19
while if x ∈ H, then
15
x is a man with probability 31 . Then p can be obtained by drawing
an appropriate tree diagram (similar to those in the solutions of the
problems in Comment No. 11 of Chapter 22). It comes out as
4 15 1 15
p = × + ×
5 19 5 31
12 3 429
= + = (10)
19 31 589
As this is not one of the options, a candidate is alerted that the paper-
setters have in mind some other meaning of the probability of correct-
ness of the blood test. But that is likely to cause confusion. Instead
of the statement about the blood test as given, it would have been far
better had the paper-setters spelt it out by saying that a person having
the disease tests positive with probability 54 and a healthy person tests
positive with probability 15 .
It is instructive to look at the difference in the two interpretations
of the probability of correctness of a test in terms of conditional prob-
ability. In this problem, if a person is known to be diseased (through
some other means perhaps), then 54 is the probability that the test is
positive. In other words, it is the conditional probability that a person
tests positive given that he/she is diseased. Effectively, this means that
you are testing the test and not the patient! This is usually done when
the test is in its experimental stages, for example, when it is claimed
as a cheap, convenient alternative to some other well established, reli-
able but invasive and expensive test. After conducting a large number
of such trials (both on diseased as well as healthy persons), a picture
emerges about the diagnostic value of the test, that is, the conditional
probability, say p∗ , that a person has a disease given that he/she tests
positive. In real life, it is this interpretation that is relevant because
now you are testing an actual patient (whose state vis-a-vis the disease
is unknown).
So, in a sense the two interpretations of the probability of the
correctness of the test are poles apart. That makes it all the more im-
perative that the question clearly tells which interpretation is intended.

18
Q.16 The number of functions f : [0, 1] −→ [0, 1] satisfying |f (x) − f (y)| =
|x − y| for all x, y in [0, 1] is
A. exactly 1
B. exactly 2
C. more than 2, but finite
D. infinite

Answer and Comments: (B). The hypothesis implies that f is one-

to-one and also continuous. Therefore by the Intermediate Value Prop-
erty, it is either strictly increasing or strictly decreasing. (See Comment
No. 6 of Chapter 16.)
If f is strictly increasing, we claim that f (x) = x for all x ∈ [0, 1].
We certainly have 0 ≤ f (0) ≤ f (x) ≤ f (1) ≤ 1. If f (x) < x, we would
get [f (0), f (x)] ⊂ [0, x] and hence, |f (0)−f (x)| = f (x)−f (0) < x−0 =
|0−x| which contradicts the hypothesis. On the other hand, if f (x) > x,
then we would get [f (x), f (1)] ⊂ [x, 1] and hence |f (x) − f (1)| < |x − 1|
a contradiction again. So we must have f (x) = x for all x ∈ [0, 1].
By a similar reasoning, if f is strictly decreasing we show that f (x) =
1 − x for all x ∈ [0, 1]. (A slicker way is to let g(x) = 1 − f (x). Then g
also satisfies the conditions of the problem and is strictly increasing. So
we can apply the part above to get that g(x) = x and hence f (x) = 1−x
for all x ∈ [0, 1]. Hence there are exactly two functions satisfying the
conditions given in the problem, viz. f (x) = x and f (x) = 1 − x.
If, instead of [0, 1], the function f is from IR to IR, and satisfies that
|f (x) − f (y)| = |x − y| for all x, y ∈ IR, then also it follows that f is
either strictly increasing or strictly decreasing. Also it can be shown
that f is differentiable in both the cases and f ′ ≡ 1 in the first case,
while f ′ (x) ≡ −1 in the second case. Therefore, f (x) = c + x for some
constant c in the first case while in the second case, f (x) = c − x for
some constant c. So, there are infinitely many such functions and these
are the only ones.
A good problem. But the reasoning involved is not trivial to write
down rigorously. So this problem is unsuitable as an MCQ. It would
be ideal as a full length question asking to characterise all functions
f : IR −→ IR satisfying |f (x) − f (y)| = |x − y| for all x, y ∈ IR.

19
Q.17 Suppose A is a 3 ×3 matrix consisting of integer entries that are chosen
at random from the set {−1000, −999, . . . , 999, 1000}. Let P be the
probability that either A2 = −I or A is diagonal, where I is the 3 × 3
identity matrix. Then
A. P < 10118 B. P = 1
1018
2 52 54
C. 10518 ≤ P ≤ 1018
D. P ≥ 1018

Answer and Comments: (A).The possibility that A2 = −I is vac-

uous because in that case we would have (det(A))2 = det(A2 ) =
det(−I) = (−1)3 = −1 which is impossible since the determinant of A
is a real number (in fact, an integer). So P is only the probability that
the matrix A is diagonal. The given set of integers has 2001 elements.
As the 9 entries of A are independent of each other, there are, in all,
(2001)9 matrices. Out of these the diagonal matrices are those where
all three diagonal elements are some integers from these 2001 integers
and the remaining six entries are 0. So, there are only (2001)3 diagonal
(2001)3 1
matrices. Hence P = 9
= . As 2001 > 103 , we have
(2001) (2001)6
P < 10118 .
An artificial combination of matrices and combinatorics.

Q.18 Let xk be real numbers such that xk ≥ k 4 + k 2 + 1 for 1 ≤ k ≤ 2018.

Denote N = 2018k=1 k. Consider the following inequalities:
P

2018
!2 2018
!
kx2k
X X
I. kxk ≤N
k=1 k=1
2018
!2 2018
!
k 2 x2k
X X
II. kxk ≤N
k=1 k=1

Then

A. both I and II are true

B. I is true and II is false
C. I is false and II is true
D. both I and II are false

20
Answer and Comments: (A). The sums on the L.H.S. in both the
inequalities are the same, while the sum on the R.H.S. of I is clearly
smaller than that of II. So, if I is true, then automatically so is II.
To check if I is true it is very tempting to apply the Cauchy-Schwarz
inequality to the L.H.S. In its classic form it says that for any positive
integer n and any real numbers u1 , u2 , . . . , un , v1 , v2 , . . . , vn we have
n
!2 n
! n
!
u2k vk2
X X X
uk vk ≤ (1)
k=1 k=1 k=1

If we apply this to the L.H.S. of I taking n = 2018, uk = k and vk = xk ,

as seems most natural to do, we get
2018
!2 2018
! 2018
!
2
x2k
X X X
kxk ≤ k (2)
k=1 k=1 k=1

The first sum on the R.H.S. equals n(n+1)(2n+1)

6
where n = 2018. Also
2018
X 2018 × 2019
N = k= . With this substitution, I will follow if we
k=1
2
can show
2018 2018
! !
2037
x2k kx2k
X X
≤ (3)
3 k=1 k=1

But this may not be true because the coefficient of x2k in the L.H.S. is
a fixed number 2037
3
while that on the R.H.S. is k for k = 1, 2, . . . , 2018.
So nearly two third terms in the sum in the L.H.S. are greater than the
corresponding terms in the sum on R.H.S. Moreover, no upper bounds
are given on any xk . So even a single huge x1 can render (3) false.
The failure of this approach can be analysed if we recall when
equality holds in the Cauchy-Schwarz inequality. In (1) it holds if
and only if the vk ’s are proportional to the uk ’s. So, in (2) equality
holds, if and only if x1 , x2 , . . . , x2018 are proportional to 1, 2, . . . , 2018
respectively. This precludes the possibility of choosing x1 very huge as
compared to the other xk ’s. So, in (2) the R.H.S. is likely to be a lot
bigger than the L.H.S. Hence it will not be of any help in proving (3).
Nothing wrong in aiming at a higher target, because if you can shoot it
then you can comfortably shoot the lower one. But the target you are

21
 aiming at should not be so high that it is beyond your capacity. So,
it was a mistake to apply the Cauchy Schwarz inequality in the most
obvious way.
Let us see if we can apply the same inequality a little more efficiently.
2018 2018
k 2 . Had the
X X
A hint is provided by giving N as k rather than as
k=1 k=1
latter been the case, then, the constant on the L.H.S. of (3) would
go away and then (3) would be indeed be true. If the first sum on the
2018
X √
R.H.S. of (2) were k then it would exactly equal N. But k = ( k)2 .
k=1 √ √
So if we rewrite kxk as k × kxk and then apply the Cauchy-Schwarz
inequality, then, instead of (2) we would get

2018
!2 2018
! 2018
! 2018
!
kx2k kx2k
X X X X
kxk ≤ k =N (4)
k=1 k=1 k=1 k=1

and now we are home because the R.H.S. of (4) is the same as that of
I. Thus we have proved I and as already observed II as well.
A somewhat tricky problem. Also the lower bounds given on xk ’s
serve no purpose. Such redundant stipulations only serve to confuse
the candidates.

Q.19 Let x2 = 4ky, k > 0, be the parabola with vertex A. Let BC be its latus
rectum. An ellipse with centre on BC touches the parabola at A, and
cuts BC at points D and E such that BD = DE = EC (B, D, E, C
in that order). The eccentricity of the ellipse is
√ √
(A.) √1 (B.) √1 (C.) 5
(D). 3
2 3 3 2

Answer and Comments: (C). The eccentricity of an ellipse depends

only on the relative proportion of the lengths of its axes. So, we may as
well take k = 1, just as in Q.5 about a rectangular hyperbola x2 − y 2 =
a2 we took a = 1 without loss of generality. Then the data of the
problem can be pictured as in the figure below.

22
 y

2
x = 4y

(−2,1) (−2/3,1) (2/3,1)

B D E C (2,1)

x
A

The ellipse has its semi-minor axis 23 and semi-major axis 1. (Note that
the major axis is along the y-axis and not parallel to the
√ x-axis. Hence
4
its eccentricity is given by 9 = 1 − e which gives e = 35 .
2

A simple problem, once the data is correctly understood. Note that

it is not necessary to write down the equation of the ellipse. All we
need is the lengths of its axes.

Q.20 Let f : [0, 1] −→ [−1, 1] and g : [−1, 1] −→ [0, 2] be two functions such
that g is injective and g ◦ f : [0, 1] −→ [0, 2] is surjective. Then
A. f must be injective but need not be surjective.
B. f must be surjective but need not be injective.
C. f must be bijective
D. f must be a constant function

Answer and Comments: (B). Call g ◦ f as h. Since h is surjective,

so is g. To see this take any z ∈ [0, 2]. Then there is some x ∈ [0, 1]
such that z = h(x) = g(f (x)) = g(y) where y = f (x) ∈ [−1, 1]. As
g is also given to be injective, we get that it is a bijection. So its
inverse g −1 : [0, 2] −→ [−1, 1] exists. From h = g ◦ f we now get
g −1 ◦ h = g −1 ◦ (g ◦ f ) = (g −1 ◦ g) ◦ f = f . So, f is the composite of
two surjective functions, h and g −1 (the latter actually being bijective).
Hence f is surjective.
The question now is whether f is injective too. We claim that
it need not be so. We fix some bijection g : [−1, 1] −→ [0, 2], say

23
g(y) = y + 1. For f to be surjective, it is enough if it maps some part,
say [0, 1/2] of [0, 1] onto [−1, 1]. The rest half, viz. (1/2, 1] can be
mapped arbitrarily to any points in [−1, 1]. The data of the problem
will be satisfied because g is a bijection. For an actual counter-example,
we construct f by taking it of the form f (x) = Ax + B for some real
numbers, chosen so that f (0) = −1 and f ( 21 ) = 1. This gives B = −1
and A2 + B = 1. So, A = 4. Hence f (x) = 4x − 1 for x ∈ [0, 21 ] and for
x ∈ ( 21 , 1] choose f (x) to be any element of [−1, 1].

24
 Section 2

Q.81 Let R be a rectangle, C a circle, and T be a triangle in the plane. The

maximum possible number of points common to the perimeters of R,
C and T is

A. 3 B. 4 C. 5 D. 6

Answer and Comments: (D).

R T

It is said that a good cartoon is one that speaks for itself. No

captions, no words in the mouths of the personalities caricatured are
needed. Such verbal aids are considered a disqualification both for the
artistic quality of the cartoonist and his/her sense of humour.
In mathematics, it is hard to lay down what constitutes a good
problem. But if we take it to mean one whose complete solution is a
diagram without a single word needed to explain it, then this problem
meets that requirement. The figure above shows a rectangle R (which
is really a square), a circle C and a triangle T and six points (marked in
colour) that are common to the perimeters of all three. Nothing more
needs to be said. Had the paper-setters given 7 (or any number higher
than 6) as a possible answer, then some reasoning would have to be
given. Not that it is difficult. Given 7 or more points on the perimeter
of a triangle T , at least three would have to lie on one of the three sides.
But then they cannot all lie on a circle. (The first part of the reasoning

25
is sheer common sense. But it has a formal name, the pigeon hole
principle. It says that if a set with more than rs elements (with r and
s some positive integers) is expressed as a union of r subsets, then at
least one of these subsets contains more than s elements.
In many applications of the pigeon hole principle, little more is
needed. The classic example of this is the Ramsey problem. Suppose
that there are six persons at a party. Prove that either there are three
among them every two of whom know each other or that there are three
no two of whom know each other. (There is a more colourful version of
this. Suppose that there are six points in plane, and every two of them
are joined by an arc. We ignore the intersections of these arcs, except
when the intersection is at one of these six points. Ramsey problem asks
to show that no matter how we colour these 15 arcs with two colours,
say red and blue, there will be a monochromatic triangle, i.e. a triangle
all three of whose edges are of the same colour.) The problem is simple
but not trivial. For a solution see Comment No. 16 of Chapter 6. But
Ramsey’s problem is just the beginning. Using generalised versions of
the pigeon hole principle, highly non-trivial results can be obtained.
There is a whole branch of mathematics, called Ramsey theory.
Of course, just because there cannot be more than 6 points common
to the perimeters of R, C and T , does not, by itself, mean that there
can be a situation where there are six points. Perhaps by a more clever
argument than the pigeon hole principle, one can get a better upper
bound. This happens frequently in mathematics. For example, for an
acute angle θ, we have sin θ cos θ ≤ 1 because both the factors are at
most 1. But if we rewrite sin θ cos θ as 21 sin 2θ, we see that it can be at
most 12 . This is the best upper bound because it is actually attained
when θ = π4 and so we cannot improve on it.
Returning to the problem, the deceptive word in it is ‘rectangle’.
Normally, we take a rectangle so that its length is considerably bigger
than its breadth. It is clear that if R is such a rectangle, then a circle
can cut it only on its longer sides. But if the rectangle is more like
a square, then it is possible to draw a circle C which cuts all its four
sides in two points each. This will happen if the circle√ has the same
centre as the square and its radius lies between a and 2a where 2a is
the side of the square.

26
 Once this idea strikes, the solution is easy. We draw such a square
R and the circle C as shown in the figure. We take both the points of
their intersection on one pair of parallel sides. But from the other two
sides we pick only one each so that they lie on a line parallel to the
other two sides. (This can be done in several ways, But we need only
one.) Then we draw a triangle T in such a way that each side contains
two of these points. One of the sides contains an edge of the square
too. So, the points chosen on it qualify to lie both on R and T .
Some persons hesitate to accept a square as a rectangle. (It all
depends on how inclusive you are about degeneracies. A square is a
degenerate case of a rectangle. But then a circle is a degenerate case
of an ellipse too.) Anyway, if instead of a square we want a ‘genuine’
rectangle, i.e. a rectangle which is not a square, that can be arranged
by elongating the two vertical sides of the square downwards slightly
and pushing the base of T downwards accordingly.

Q.82 The number of different possible values for the sum x + y + z, where
x, y, z are real numbers such that x4 + 4y 4 + 16z 4 + 64 = 32xyz is

A. 1 B. 2 C. 4 D. 8

Answer and Comments: (C). Normally a single equation in more

than one variables has infinitely many solutions. A single equation in
three variables x, y, z generally describes a surface such as a plane, or
a sphere or an ellipsoid in the three dimensional space IR3 . When that
happens, x + y + z can take infinitely many distinct values, the only
exception being when the surface described by the equation is itself a
plane of the form x + y + z = k for some constant k.
So, obviously, the present problem is not about identifying the
solutions of a generic equation in x, y, z. There has to be something
very special about that equation. Already we encountered one problem
of a similar spirit in Q.6. There an optimisation problem was converted
to a problem of solving an equation. In the present problem too, if we
write every term on the L.H.S. as a fourth power, the given equation
becomes
√ √
x4 + ( 2y)4 + (2z)4 + (2 2)4 = 32xyz (1)

27
 If x, y, z are all non-negative,
√ then by √the A.M.-G.M. inequality the
L.H.S. is at least 4 × x × 2y × 2z × 2 2. But this product is simply
32xyz which is the R.H.S. of (1). So, the only possible solution is
√ √
x = 2y = 2z = 2 2 (2)

i.e.
√ √
x = 2 2, y = 2, z = 2 (3)

So the given equation has only one solution if x, y, z are all positive.
But if we change the signs of any two values in (3), the equation still
holds. This can be done in three different ways. So, in all the given
equation has four distinct solutions.
In commenting on Q.6 we already remarked about the duplication
of ideas. We now have a case of triplication!

Q.83 Let Γ be a circle with diameter AB and centre O. Let ℓ be the tangent
to Γ at B. For each point M on Γ different from A, consider the tangent
t at M and let it intersect ℓ at P . Draw a line parallel to AB through
P intersecting OM at Q. The locus of Q as M varies over C is

A. an arc of a circle
B. a parabola
C. an arc of an ellipse
D. a branch of a hyperbola

Answer and Comments: (B). We first give a pure geometry solution.

(For notational simplicity, the symbols Γ, ℓ and t are replaced by C, L
and T respectively in the figure.)
Through O draw a line parallel to ℓ and let it cut P Q at S.

28
 A

Q
M
C

S
O
T

L
B P

The triangles OSQ and P MQ are right angled and have one angle
common. Also OS = BP = P M. Hence ∆OSQ ≡ ∆P MQ. Therefore
OQ = QP . Hence the point Q is equidistant from the fixed point O
and the fixed line ℓ. Therefore its locus is a parabola with focus O and
directrix ℓ.

A solution based on coordinates is almost a copy of this. Take O as

the origin and the axes parallel to ℓ and AB. Without loss of generality
assume that Γ has radius 1. Then its equation is x2 + y 2 = 1 and that
of ℓ is y = −1. The point M can be taken as M = (cos θ, sin θ) for
some θ 6= π2 . The equation of the tangent at M is
x cos θ + y sin θ = 1 (1)
Solving this with y = −1 gives P as ( 1+sin
cos θ
θ
, −1) = (tan θ + sec θ, −1).
The equation of OM is y = tan θ. This gives
Q = (tan θ + sec θ, tan θ(tan θ + sec θ) (2)
So, the equation of the locus of Q is obtained by eliminating θ from
the two equations
x = tan θ + sec θ (3)
and y = tan θ(tan θ + sec θ) (4)
The elimination can be done by squaring x − tan θ = sec θ and putting
tan θ = xy . This gives
y y2
(x − )2 = 1 + 2 (5)
x x
29
 which simplifies to (x2 − y)2 = x2 + y 2 and further to

x2 − 2y = 1 (6)

which is a parabola.
Clearly, the pure geometry solution is shorter and gives more specific
information about the locus, viz. its focus and directrix. Of course, we
can get this from (6) too.

Q.84 The number of solutions x of the equation sin(x + x2 ) − sin(x2 ) = sin x

in the interval [2, 3] is

A. 0 B. 1 C. 2 D. 3

Answer and Comments: (C). We begin by rewriting the given equa-

tion as
x x
2 cos( + x2 ) sin = sin x (1)
2 2
Rewriting the R.H.S. as 2 sin x2 cos x2 , gives two possibilities, (i) sin x2 = 0
and (ii) cos( x2 + x2 ) = cos x2 .
The possible solutions of (i) are x = 2nπ for some integer n. This is
not possible for any value of x in [2, 3].
However, (ii) holds whenever x2 +x2 = 2nπ± x2 for some integer n. The
2
√ x = 2nπ which is possible for x ∈ [2, 3]
choice of the + sign gives only
when n = 1 and x = 2π. The choice of the negative sign gives x2 +x =
2nπ which again is possible for only one value of x ∈ [2, 3]. (It is not
necessary to identify this value. It is enough to note that the function
f (x) is strictly increasing on [2, 3] and therefore has range [6, 12]. There
is only one multiple of 2π, viz. 2π in [6, 12].) Put together, the given
equation has two solutions for x ∈ [2, 3]. (Throughout, we need that π
lies between 3 and 4.)
A very simple problem, hardly belongs to Section 2.

Q.85 The number of polynomials p : IR −→ IR satisfying p(0) = 0, p(x) > x2

for all x 6= 0, and p′′ (0) = 12 is

A. 0

30
 B. 1
C. more than 1, but finite
D. infinite

Answer and Comments: (A). We can get rid of x2 by considering

a new polynomial q(x) = p(x) − x2 . Although p and q determine
each other uniquely, it is a little easier to work with q(x) because the
conditions on q(x) are q(0) = 0, q(x) > 0 for all x 6= 0 and q ′′ (0) =
1
2
− 2 = − 32 .
Assume that q(x) = a0 + a1 x + a2 x2 + a3 x3 + . . . + an xn for some
integer n ≥ 0 and some real numbers a0 , a1 , . . . , an . The requirement
that q(0) = 0 means a0 = 0. We claim that a1 also vanishes. Note that
q ′ (0) = a1 . So, if a1 6= 0, then depending upon the sign of a1 , there
will be points on one of the two sides of 0 where q(x) < q(0). But this
contradicts the requirement that q(x) > 0 for all x 6= 0.
So, q(x) is of the form a2 x2 +a3 x3 +. . .+an xn . But then q ′′ (0) = 2a2 .
The requirement that q ′′ (0) = − 32 So we get a2 = − 34 .
Summing up, q(x) is of the form − 43 x2 + a3 x3 + . . . + an xn . Since there
are no terms of degree 0 and 1, the dominating term near 0 is − 34 x2
and it is negative for non-zero x sufficiently close to 0. For a rigorous
proof of this, write q(x) as
3
q(x) = − x2 (1 + r(x)) (1)
4
where

r(x) = b3 x + b4 x2 + . . . + bn xn−2 (2)

ai
and bi = −3/4 for i = 3, 4, . . . , n. Then r(x) → 0 as x → 0 and therefore
we can find some δ > 0 such that − 12 < r(x) < 12 for all x ∈ (−δ, δ). But
then, 12 < 1 + r(x) < 23 for all x ∈ (−δ, δ). Hence for x ∈ (−δ, δ) − {0},
1 + r(x) > 0 and so from (1), q(x) < 0, a contradiction.
Therefore there are no polynomials satisfying the given conditions.
A very good problem, requiring only that the sign of a polynomial
near 0 is determined by the sign of its smallest non-vanishing coefficient.

31
Q.86 Suppose the limit
√ Z 1 1
L = lim n dx
n→∞ 0 (1 + x2 )n
exists and is larger than 12 . Then
1
A. 2
<L<2 B. 2 < L < 3 C. 3 < L < 4 D. L ≥ 4

Answer and Comments: (A). The problem does not ask us to iden-
tify L or even to prove that it exists. It is also given that L > 12 . Our
task is to decide which of the four possibilities hold for L. Call the
given integral as In . If we could evaluate In directly as√a function of
n, we can decide which possibility holds by considering nIn for large
values of In .
The substitution x = tan θ gives
π/4 π
Z
I1 = 1dθ = (1)
0 4
However, for n > 1, it is not easy to express In directly as a function of
n. The best we can do is to get a reduction formula for it, i.e. express
In in terms of In−1 . This can be done if we integrate by parts.
Z 1
In = (1 + x2 )−n dx
0
1 Z 1
= [x(1 + x2 )−n ] + 2n (1 + x2 )−n−1 x2 dx
0 0
1 Z 1
= n + 2n [(1 + x2 )−n−1 (1 + x2 − 1)dx
2 0
1
Z 1 Z 1
2 −n
= n + 2n (1 + x ) dx − 2n (1 + x2 )−n−1 dx
2 0 0
1
= n + 2nIn − 2nIn+1 (2)
2
Hence
1 2n − 1
In+1 = + In (3)
2n+1 n 2n
By induction on n we shall show that
√
nIn < 2 (4)

32
for all n. This does not quite prove that L < 2. But it does prove that
L ≤ 2 and eliminates options (B), (C) and (D). Surely (4) is true for
n = 1 since I1 = π4 < 2.
For the inductive step we use (3) and the inductive hypothesis to get
√ √
√ n + 1 2n − 1 n + 1 √
n + 1In+1 = + √ nIn
2n+1 n 2n n
√ s
n + 1 2n − 1 1
< + 1 + (5)
2n+1 n n n
√
So n + 1In+1 < 2 will follow if we can prove that
s √
2n − 1 1 n+1
1 + < 2 − n+1 (6)
n n 2 n
√
The substitution h = n1 reduces the L.H.S. to (2 − h) 1 + h. By the
Taylor approximation of order 2, we have
√ h h2
1+h<1+ − (7)
2 8
for all h > 0. Hence
√ h h2
(2 − h) 1 + h < (2 − h)(1 + − )
2 8
2 3
3h h
= 2− +
4 8
3 1
= 2− 2 + 3 (8)
4n 8n
Hence (6) will be proved if we can show that
√
3 1 n+1
2 − 2 + 3 < 2 − n+1 (9)
4n 8n 2 n
which can be rewritten as
√
n+1 3 1 6n − 1
n+1
< − 2 = (10)
2 4n 8n 8n2

33
or equivalently,
√
n + 1n2
< 2n−2 (11)
6n − 1
The L.H.S. is bounded by a rational function in n while the R.H.S. is
an an exponential function. Hence (11) holds for all sufficiently large n.
But
√ by a direct verification, it is true even for n = 2, 3, 4, . . . . Therefore
nIn < 2 for all n ≥ 2. (For the sake of completeness, an elementary
inductive proof of (11) can be given by observing that as n increases
by 1, the expression on the L.H.S. increases by a factor which is only
marginally bigger than 1, while the R.H.S. always increases by a factor
of 2.)
This proof is far too advanced for the KVPY standards. A much
shorter and more direct proof of (4) (which does not use the reduction
formula (3)) can be given simply by obtaining an upper bound on the
1
integrand on [0, 1]. All we need is to expand (1 + x2 )n by the
(1 + x2 )n
binomial theorem. As all the terms are positive, we get
(1 + x2 )n > 1 + nx2 (12)
for all n > 1. Hence
1 dx 1 dx
Z Z
In = 2 n
< (13)
0 (1 + x ) 0 1 + nx2
√
The substitution y = nx converts the integral on the R.H.S. to
1 1 dy
Z
√ π
√ which we already know to be n 4 . Hence,
n 0 1 + y2
√ π
nIn < (14)
4
for all n ∈ IN. This implies that L ≤ π4 and hence L < 2 since π4 < 2.
A reduction formula is a powerful tool. But when the question is
not so demanding, an elementary alternative can often be found as in
this problem.
We indicate how a similar reduction formula for a closely related
integral, along with a well-known approximation formula for the fac-
torials (called the Stirling approximation) enables us to obtain the
value of L. (Of course, this is not asked in the present problem.)

34
Putting x = tan θ, we get
Z π/4
In = cos2n−2 θdθ (15)
0
= Kn − J n (16)

where
Z π/2
Kn = cos2n−2 θdθ (17)
0
Z π/2
and Jn = cos2n−2 θdθ (18)
π/4

π 1 √
Note that Jn ≤ 4 2n−1
, nJn → 0 as n → ∞. Hence
So
√ √
lim nIn = lim nKn (19)
n→∞ n→∞

provided the second limit exists.

There is considerable advantage in replacing In by Kn . Like In ,
there is a reduction formula for Kn too. But it is a lot simpler. We
first note that
Z π/2
2n−2 π Z π/2
Kn = cos ( − θ)dθ = sin2n−2 θdθ (20)
0 2 0

As in the derivation of (2), we integrate by parts to get

Z π/2
Kn = sin2n−2 θdθ
0
π/2 Z π/2
= [− sin2n−3 θ cos θ] + (2n − 3) sin2n−4 θ cos2 θdθ(21)
0 0
Z π/2
= (2n − 3) sin2n−4 θ(1 − sin2 θ)dθ
0
= (2n − 3)(Kn−1 − Kn ) (22)

which implies
2n − 3
Kn = Kn−1 (23)
2n − 2
We could have also derived (2) using the substitution x = tan θ. But
the advantage this time is that the first term in (21) vanishes and so

35
the reduction formula is much simpler. Indeed, repeated applications
of this formula give
2n − 3
Kn = Kn−1
2n − 2
(2n − 3)(2n − 5)
= Kn−2
(2n − 2)(2n − 4)
..
.
(2n − 3)(2n − 5) . . . 1
= . . . K1 (24)
(2n − 2)(2n − 4) . . . 2
Z π/2 π
By a direct calculation, K1 = 1dθ = . Multiplying and dividing
0 2
by (2n − 2)(2n − 4) . . . 2 we get

(2n − 2)! π
Kn = (25)
[2n−1 (n− 1)!]2 2

To go from this point onwards, we need the following result, called

Stirling approximation. Its proof can be given by approximating
the area under the curve y = ln x, 1 ≤ x ≤ n by a collection of trapezea
with bases the intervals [k, k + 1], k = 1, 2, . . . , n − 1 on the x-axis. But
we shall not go into it except to remark that the role of the function
n
X
f (x) = ln x is not surprising because ln(n!) = f (k).
k=1

n! √
Theorem: As n → ∞, √ n n → 2π
n( e )
nn√
This theorem allows us to replace n! by for large n. Doing
2πn
en
so for the two factorials (2n − 2)! and (n − 1)! in (25), we get
√
√ π n(2n − 2)!
lim nKn = lim n−1
n→∞ 2 n→∞ [2 (n − 1)!]2
q √ 
π 2π(2n − 2)(2n − 2)2n−2 n e2n−2
= lim  × 2n−2 × 
2 n→∞ e(2n−2) 2 2π(n − 1)2n−1

36
 √ √
The bracketed expression simplifies to √1π × n−1
n−1
n
and tends to √1π
√ √ √
as n
√
→ ∞. Hence nKn → 2π as n → ∞, By (19) √nIn also tends
to 2π as n → ∞. In other words, L exists and equals 2π .

Q.87 Consider the set An of points (x, y) such that 0 ≤ x ≤ n, 0 ≤ y ≤ n

where n, x, y are integers. Let Sn be the set of all lines passing through
at least two distinct points of An . Suppose we choose a line ℓ at random
from Sn . Let  Pn be the probability
 that ℓ is tangent to the circle
2 2 2 1 2
x + y = n 1 + (1 − n ) . Then lim Pn is
√
n→∞
√
A. 0 B. 1 C. 1/π D. 1/ 2

Answer and Comments : (A). The statement of the problem is long-

winded. The set An has (n + 1)2 = n2 + 2n + 1 elements. From these
we can choose two distinct points in 21 (n2 + 2n + 1)(n2 + 2n) ways. But
that does not mean that in the set Sn , there are these many distinct
lines, because the same line may contain more than two elements of
An . Counting how many lines there are in Sn is not an easy task. Even
if we are able to do that, we further have to count how many of these
lines touch the circle, say C, given by
!
1
C = {(x, y) : x2 + y 2 = n2 1 + (1 − √ )2 (1)
n

So, calculating Pn as an explicit function of n is a formidable task.

Fortunately, the problem does not ask that. All it asks is the limit,
lim Pn .
n→∞
And there is a way to see that this limit is 0, by getting an upper
bound on Pn . The key idea is that if a line ℓ touches a circle then the
points in An , if any, that lie on it cannot be inside the circle. In other
words all such points, say P or Q, lie in the portion, say R, of the
square [0, n] × [0, n] outside the arc AB of the circle C shown in the
figure below (where ℓ is shown as L).

37
 y

D
C

(0,n) E(n,n)
B Q R
P
A
L

x
(0,0) (n,0)

Consider the circle D centred at (0, 0) and passing through E =

(n, n). Its equation is

D = {(x, y) : x2 + y 2 = 2n2 } (2)

Then the region R lies in the first quadrant in between the concentric
circles C and D. So, if we denote the area of R by ∆n , we get an easy
upper bound on it, viz.,

πn2
" !#
1
∆n ≤ 2 − 1 + (1 − √ )2
4 n
2
!
πn 2 1
= √ −
4 n n
πn2
≤ √ (3)
2 n

We can get a much better upper bound on ∆n . In fact, we can even

calculate it exactly. But the upper bound in (3) is sufficient for our
purpose. The points of the set An are evenly distributed in the square
[0, n] × [0, n] whose area is n2 . So, the probability that a point taken at
random from this square belongs to R is approximately the ratio of the
area of R to that of the square. Because of the upper bound above this
π
probability is at most √ . And the probability that two such points
2 n
both lie in R will be even smaller. So, without calculating Pn exactly,

38
 π
we know that it is at most √ . As this tends to 0 when n → ∞, so
2 n
does Pn .
An excellent problem which tests a candidate’s ability to get to the
key idea, which in the present problem is that if ℓ touches the given
circle C then no point on it can be inside C. After that there is a jump
in concluding that the probability that a point be picked randomly from
some subset R of the square [0, n] × [0, n] is proportional to its area. It
will be easier to accept this if instead of the set An as given we consider
the set Bn of all ordered pairs of the form ( nx , ny ) where x, y are integers
ranging from 0 to n. In effect, Bn is obtained by compressing the square
[0, n] × [0, n] to the unit square S = [0, 1] × [0, 1] under the mapping
which sends a point (x, y) to the point ( nx , ny ). This compression does
not affect tangency and so the probability Pn remains the same. Note
that S is independent of n. But as n increases, not only the number
of points in Bn increases, their density also increases. As n → ∞, the
discrete set Bn of points tends to the continuous set of all points of the
unit square S. The concept of probability as the ratio of cardinalities
no longer applies now. Instead we have to take the ratios of their areas.
Problems of this type are not asked commonly. To help understand
the solution better, at the end of Section 9 of Chapter 23 on infinitistic
probability, there is a problem of a similar type. It asks to find lim pn
n→∞
where pn is the probability that two points picked at random from
the set { 2i
n
: 0 ≤ i ≤ n} are at a distance 1 or less apart from each
other. Here pn can be calculated explicitly and then its limit comes
out as 34 . But if our interest is only in the limit, it is the same as the
probability that two points taken at random from the interval [0, 2] are
at a distance 1 or less apart. And this is solved at the beginning of the
same section, by taking the ratio of certain areas.

Q.88 Let f : [0, 1] −→ IR be an injective continuous function that satisfies

the condition
−1 < f (0) < f (1) < 1.
Then the number of functions g : [−1, 1] −→ [0, 1] such that (g◦f )(x) =
x for all x ∈ [0, 1] is

A. 0

39
 B. 1
C. more than 1, but finite
D. infinite

Answer and Comments: (D). The hypothesis implies that f is either

strictly increasing or strictly decreasing on [0, 1]. (This is a consequence
of the Intermediate Value Property for continuous functions and was
also used in the solution to Q.16.) Since we are given that f (0) < f (1),
only the first possibility holds. Moreover, the interval [f (0), f (1)] is
the range of the function f . This is given to be a proper subset of
the interval [−1, 1]. f maps [0, 1] bijectively onto [f (0), f (1)] which is
a proper subset of [−1, 1] which is the domain of g. For every y ∈
[f (0), f (1)], there is a unique x ∈ [0, 1] such that f (x) = y. So the
requirement (g ◦ f )(x) = x implies g(y) has to be x. But when y ∈
[−1, 1] − [f (0), f (1)], we are free to define g(y) to be any element of
[0, 1], without violating the data. So there are infinitely many choices
for g.
There is little point in asking this question, especially in Section 2,
after having asked a much better question (Q.16) in Section 1.
Q.89 The maximum possible area bounded by the parabola y = x2 + x + 10
and a chord of the parabola of length 1 is
1 1 1 1
A. 12
B. 6
C. 3
D. 2

Anwer and Comments: (B). We can recast the equation of the

parabola as y − 39 4
= (x + 21 )2 . So, if we shift the origin to the point
(− 21 , 39
4
) this will reduce to the standard form Y = X 2 where X =
x + 21 , Y = y − 394
. But the concepts of length and area are independent
of a particular frame of reference. So, we might as well suppose that
the equation of the parabola is
y = x2 (1)
(An alert candidate can save time by getting this directly by realising
that the particular point to which the origin is shifted is unimportant.)
Let AB be any chord of it, where A = (a, a2 ) and B = (b, b2 ) and for
some a, b ∈ IR and we assume, without loss of generality that a < b.

40
 y

A ∆

x
a O h b

This is a vertically upward parabola which always lies on or above

the x-axis. Moreover , as the function f (x) is concave upward, the
chord AB lies above the arc AB, the two meeting only at A and B.
So, the area, say ∆ bounded by the arc AB and the chord AB is

∆ = ∆1 − ∆2 (2)

where ∆1 is the area between the chord and the x-axis and ∆2 is the
area between the curve y = x2 and the x-axis, both lying between x = a
and x = b.
As ∆1 is the area of a trapezeum, a direct calulation gives

a2 + b2 Z b
∆ = (b − a) − x2 dx
2 a
a2 + b2 b3 − a3
= (b − a) − (3)
2 3
As the factor b − a is common to both the terms, it is conveninet to
call it h and write b as a + h. Then (3) becomes

2a2 + 2ah + h2 3a2 h + 3ah2 + h3

∆ = h−
2 3
1 2 h3
= (6a h + 6ah2 + 3h3 − 6a2 h − 6ah2 − 2h3 ) = (4)
6 6
Note also that h is the vertical projection of the chord AB onto the
x-axis. As the chord is given to be of length 1, we have h ≤ 1 with the

41
 maximum occurring when h = 1, i.e. when the chord is horizontal. So
the maximum value of ∆ is 61 .
The problem is conceptually simple. But the introduction of h
makes it computationally simple too because the expression for ∆ in
terms of h is surprisingly simple. (It is also remarkable that ∆ depends
only on b − a and not on both a and b.) Had we not introduced h,
instead of (4) we would get
(b − a)3
∆= (5)
6
and the problem would have been to maximise (b − a)3 subject to the
constraint (b − a)2 + (b2 − a2 )2 = 1 (obtained by setting the length of
the chord AB equal to 1). If we let θ be the angle the chord AB makes
with the x-axis, then b2 − a2 = (b − a) tan θ and the constraint becomes
b − a = cos θ and we have to maximisie cos3 θ. Clearly, the maximum
occurs when cos θ = 1. This gives the same answer as before.
Parabolas, like other conics, have been extensively studied in classi-
cal geometry. But the formula (4) seems not to be well-known, despite
its simplicity. A possible explanation is that it deals with areas unlike
most other formulas which deal mostly with angles and lengths.
√
Q.90 Suppose z is any root of 11z 8 + 20iz 7 + 10iz − 22 = 0, where i = −1.
Then S = |z|2 + |z| + 1 satisfies
A. S ≤ 3 B. 3 < S < 7
C. 7 ≤ S < 13 D. S ≥ 13

Answer and Comments: (B). As an elementary simplification, we

call iz as w. Since i8 = 1 and i7 = −i, the given equation becomes
11w 8 − 20w 7 + 10w − 22 = 0 (1)
which is equivalent to the original equation since z and w determine
each other. Further, since |w| = |z|, S = |w|2 + |w| + 1 and so the four
given options remain the same. Clearly the value of |w| determines S
and the four options correspond to (A) |w| ≤ 1, (B) 1 < |w| < 2, (C)
2 ≤ |w| < 3 and (D) |w| ≥ 3.
So, the problem asks to identify which of these four subsets of
the complex plane contains all the roots of (1). Unfortunately, it is

42
impossible to answer this by actually identifying all the roots of (1).
But we can recast it as

11w 8 + 10w = 20w 7 + 22 (2)

and hence further as

w(11w 7 + 10) = 2(10w 7 + 11) (3)

The advantage of this recasting is not that we can solve the equation,
but that when we take the squares of the absolute values of both the
sides, some of the terms in the expansion of |11w 7 +10|2 will be common
to some of the terms in the expansion of |10w 7 + 11|2. Doing so, we get
from (3),

|w|2|11w 7 + 10|2 = 4|10w 7 + 11|2 (4)

Both the sides have two factors both of which are positive. This does
not mean that their factors match. But it does imply that when one of
the factors on one side is greater than the corresponding factor on the
other side, the remaining two factors on the two sides must satisfy the
opposite inequality. That is,

|w| ≥ 2 ⇔ |11w 7 + 10|2 ≤ |10w 7 + 11|2 (5)

i.e.

|w| ≥ 2 ⇔ 121|w|14 + 110w 7 + 110w7 + 100

≤ 100|w|14 + 110w 7 + 110w7 + 121
⇔ 21|w|14 ≤ 21
⇔ |w| ≤ 1 (6)

Summing up, |w| ≥ 2 holds if and only if |w| ≤ 1 holds. But these two
statements are mutually contradictory and can never hold together. So
both are false. This gives |w| < 2 and |w| > 1. Thus 1 < |w| < 2 and
as already noted, S = |w|2 + |w| + 1 lies between 3 and 7. So (B) holds.
Lest this solution appears too tricky, it is instructive to recast (3)
a little differently, viz.
w 10w 7 + 11
= (7)
2 11w 7 + 10

43
Let us put z = w 7 . (Earlier we let w = iz and hence z = −iw. But
that is all over now. So no harm in using z for something else now.
Those who are uncomfortable, may denote w 7 by the complex number
Z instead of z.) Then (7) can be written as
w
= T (z) (8)
2
where
az + b
T (z) = (9)
cz + d
with a = d = 10 and b = c = 11. T (z) is a ratio of two linear expressions
(i.e. two first degree polynomials) in z. For this reason, T (z) is an
example of what is called a fractional linear transformation or
FLT for short. It is also called a Mobius transformation. More
az+b
generally, an FLT is any transformation of the form T (z) = cz+d where
a, b, c, d are complex numbers with ad − bc 6= 0. The last stipulation is
added to exclude the degenerate case where T (z) is a constant mapping.
Note that T (z) is undefined when the denominator vanishes, i.e. when
z = − dc . It is customary to add one point ∞ to the complex plane C |
d a
and say that T (− c ) = ∞. To make up, we set T (∞) = c with the
understanding that if c = ∞, then T (∞) = ∞. (Note that a and c
cannot both vanish because of the condition ad − bc 6= 0.)
With this understanding, every FLT can be thought of as a transfor-
mation from the extended complex plane C | ∗ = C
| ∪ {∞} to itself.

It is easy to show that it is always a bijection. The study of FLT’s is

important because all the familiar transformations of the plane such as
rotations, dilations, contractions, translations can be recognised as suit-
able FLT’s. But the most interesting are FLT’s of the type T (z) = 1z .
This transformation maps the unit circle |z| = 1 onto itself, but its
inside, i.e. the set {z ∈ C| ∗ : |z| < 1} to its outside, viz. the set
∗
{z ∈ C | : |z| > 1} and vice versa. Because of this peculiar property,
it is called an inversion. Inversions, combined with other transfor-
mations can be used to map bounded regions to unbounded ones and
vice versa and this is one of the reasons the Mobius transformations
are important in applications.
Rich as it is, we shall not go into the theory of Mobius transformations
here. The only property relevant to us is proved below.

44
Theorem: Suppose that a, b, c, d are real numbers with a = d, b = c
and |a| < |b|. Then the transformation (9) maps the unit circle onto
∗
itself, the region {z ∈ C
| : |z| < 1} onto the region {z ∈ C
| : |z| > 1}
and vice versa.

Proof: The simplest case, where a = d = 0 and b = c = 1 was

already discussed above. The proof in the general case is essentially a
duplication of the argument in the solution to the problem.
(i) First assume that |z| = 1. Then proving that |T (z)| = 1 is equivalent
to proving that |az + b|2 = |bz + a|2 . As a, b are both real, the L.H.S.
equals a2 |z|2 +ab(z+z)+b2 while the R.H.S. equals b2 |z|2 +ab(z+z)+a2 .
Since |z| = 1, both the expressions are equal. (Here we do not need
that |a| < |b|.)
(ii) Now assume that |z| < 1. Then proving that |T (z)| > 1 is equiva-
lent to proving that |az + b|2 − |bz + a|2 > 0. Expanding as in (i), and
factorising, this expression is (a2 − b2 )(|z|2 − 1) which is positive since
both the factors are negative.
(iii) Assume that |z| > 1. Here the argument is similar to (ii).

10z + 11
We now go back to (7). The R.H.S. is T (w 7) where T (z) = .
11z + 10
Since |w 7| = |w|7, the conditions |z| = 1, |z| < 1 and |z| > 1 are
equivalent, respectively, to the conditions |w| = 1, |w| < 1 and |w| > 1.
On the other hand | w2 | = |w| 2
. Hence |w| > 2 holds if and only if
|T (w)| > 1 holds and by the theorem above this happens if and only if
|w| < 1.
Thus the concept of a fractional linear transformation enables us to
cast the solution more systematically. It is possible that the problem
was inspired by the property of FLT’s given in the theorem above. If
so, the paper-setters must be commended for making an imaginative
use of FLT’s. But is not clear what is gained by introducing iz instead
of z. The problem could as well have given that a complex number
z satisfies the equation 11z 8 − 20z 7 + 10z − 22 = 0. Replacing z by
iz gives a nasty twist to the problem without enhancing its academic
quality.

45
 Concluding Remarks

The Kishore Vaigyanik Protsahan Yojana (KVPY) examinations are not

gateways to any lucrative careers. So, one expects higher academic standards
from them. Problems should be thought provoking rather than computa-
tional. In this respect, the present paper is a mixed bag. All the problems
on conics can be done by routine computations.
On the other hand, Q.86 and 87 involve the asymptotic estimations, the
former of a certain integral and the other of a certain probability. Both
are very good problems. Q.81 is a problem whose solution only requires a
wordless diagram. In this respect the problem is unusual. Q.85 is also a
good problem requiring very little computation, once the idea strikes that
near 0, the sign of a polynomial is determined by its lowest degree term.
Q.90 would have been more suitable had the syllabus included fractional
linear transformations. Without them, the solution is likely to appear tricky.
The remaining five problems of Section 2 could as well have been placed
in Section 1. On the other hand, Section 1 which is supposedly for mediocre
questions contains a few good problems, e.g. Q.2 about roots of unity, Q.8
about algebraic numbers, Q.12 about evaluation of a certain integral, Q.13
about solutions of differential equations on disconnected domains, and Q.16
about isometries of the unit interval. These five problems could have been
shifted to Section 2.
The probability problem (Q.15) is appealing because of its real life set-
ting. But because of the two possible interpretations of the probability of
correctness of a blood test, it is confusing. Q.18 involves a tricky use of the
Cauchy Schwarz inequality and is a good question. But the lower bounds
given on the numbers in it serve no purpose except to create confusion.
There is considerable duplication (and, in some cases even a triplication)
of ideas. Thus both Q.4 and Q.6 are optimisation problems converted as
problems of solving equations and so is Q.83 which comes directly from the
A.M.-G.M. inequality. Two problems (Q.7 and Q.19) ask to calculate ec-
centricity of an ellipse. Some of these problems could have been dropped
to make room for some problems in combinatorics, number theory, binomial
identities, vectors and differential equations. Some of these topics may not be
explicitly included in the syllabus. But many interesting problems in num-
ber theory require little more than common knowledge of some properties of
primes and divisibility.

46
 Shockingly absent are real life problems which can be reduced to simple
algebra or sometimes even to arithmetic. Such problems can even be tried by
laymen and serve to popularise the test from which they come. Two notable
problems of this type asked in KVPY in recent years come to mind. One
was about the loss of weight in watermelons and the other about the motion
of ants on a line.

47